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230271 | https://en.wikipedia.org/wiki/Yak | Yak | The yak (Bos grunniens), also known as the Tartary ox, grunting ox, hairy cattle, or domestic yak, is a species of long-haired domesticated cattle found throughout the Himalayan region of Gilgit-Baltistan (Kashmir, Pakistan), Nepal, Sikkim (India), the Tibetan Plateau (China), Tajikistan Pamir mountains Afghanistan and as far north as Mongolia and Siberia. It is descended from the wild yak (Bos mutus).
Etymology
The English word yak originates from the . In Tibetan and Balti it refers only to the male of the species, the female being called , or in Tibetan and in Balti. In English, as in most other languages that have borrowed the word, yak is usually used for both sexes, with bull or cow referring to each sex separately.
Taxonomy
Belonging to the genus Bos, Yaks are related to cattle (Bos primigenius). Mitochondrial DNA analyses to determine the evolutionary history of yaks have been inconclusive.
The yak may have diverged from cattle at any point between one and five million years ago, and there is some suggestion that it may be more closely related to bison than to the other members of its designated genus. Apparent close fossil relatives of the yak, such as Bos baikalensis, have been found in eastern Russia, suggesting a possible route by which yak-like ancestors of the modern American bison could have entered the Americas.
The species was originally designated as Bos grunniens ("grunting ox") by Linnaeus in 1766. Still, this name is now generally considered to refer only to the domesticated form of the animal, with Bos mutus ("mute ox") being the preferred name for the wild species. Although some authors still consider the wild yak to be a subspecies, Bos grunniens mutus, the ICZN made an official ruling in 2003 permitting the use of the name Bos mutus for wild yaks, and this is now the more common usage.
There are no recognised subspecies of yak except where the wild yak is considered a subspecies of Bos grunniens.
Physical characteristics
Yaks are heavily built animals with bulky frames, sturdy legs, rounded, cloven hooves, and extremely dense, long fur hanging lower than the belly. While wild yaks are generally dark, blackish to brown in colouration, domestic yaks can be quite variable, often having rusty brown and cream patches. They have small ears and broad foreheads, with smooth horns that are generally dark in colour. In males (bulls), the horns sweep out from the sides of the head and then curve backwards; they typically range from in length.
The horns of females (cows) are smaller, at in length, and have a more upright shape. Both sexes have a short neck with a pronounced hump over the shoulders, although this is larger and more visible in males. Males weigh , females weigh . Wild (feral) yaks can be substantially heavier, bulls reaching weights of up to . Depending on the breed, domestic yak males are high at the withers, while females are high at the withers.
Both sexes have long, shaggy hair with a dense woolly undercoat over the chest, flanks, and thighs to insulate them from the cold. Especially in bulls, this may form a long "skirt" that can reach the ground. The tail is long and horselike rather than tufted like the tails of cattle or bison. Domesticated yaks have a wide range of coat colours, with some individuals being white, grey, brown, roan or piebald. The udder in females and the scrotum in males are small and hairy as protection against the cold. Females have four teats.
Yaks are not known to produce the characteristic lowing (mooing) sound of cattle, but both wild and domestic yaks grunt and squeak, which inspired the scientific name of the domestic yak variant, Bos grunniens (grunting bull). Nikolay Przhevalsky named the wild variant Bos mutus (silent bull), believing that it did not make a sound at all, but it does.
Physiology
Yak physiology is well adapted to high altitudes, having larger lungs and heart than cattle found at lower altitudes, as well as greater capacity for transporting oxygen through their blood, due to the persistence of foetal haemoglobin throughout life. Conversely, yaks have trouble thriving at lower altitudes, and are prone to suffering from heat exhaustion above about . Further adaptations to the cold include a thick layer of subcutaneous fat and an almost complete lack of functional sweat glands.
Compared with domestic cattle, the rumen of yaks is unusually large, relative to the omasum. This likely allows them to consume greater quantities of low-quality food at a time, and to ferment it longer to extract more nutrients. Yak consume the equivalent of 1% of their body weight daily while cattle require 3% to maintain condition. They are grazing herbivores, with their wild ancestors feeding primarily on grass and sedges, with some herbs and dwarf shrubs.
Reproduction and life history
Yaks mate in the summer, typically between July and September, depending on the local environment. For the remainder of the year, many bulls wander in small bachelor groups away from the large herds. Still, as the rut approaches, they become aggressive and regularly fight with each other to establish dominance. In addition to non-violent threat displays, bellowing, and scraping the ground with their horns, bull yaks compete more directly, repeatedly charging at each other with heads lowered or sparring with their horns. Like bison, but unlike cattle, males wallow in dry soil during the rut, often while scent-marking with urine or dung. Females enter oestrus up to four times a year, and females are receptive only for a few hours in each cycle.
Gestation lasts between 257 and 270 days, so that the young are born between May and June, and results in the birth of a single calf. The cow finds a secluded spot to give birth, but the calf can walk within about ten minutes of birth, and the pair soon rejoin the herd. Females of both the wild and domestic forms typically give birth only once every other year, although more frequent births are possible if the food supply is good.
Calves are weaned at one year and become independent shortly thereafter. Wild calves are initially brown in color and only later develop darker adult hair. Females generally give birth for the first time at three or four years of age, and reach their peak reproductive fitness at around six years. Yaks may live for more than twenty years in domestication or captivity, although it is likely that this may be somewhat shorter in the wild.
Husbandry
For thousands of years, domesticated yaks have been kept in Mongolia and Tibet, primarily for their milk, fibre (wool), and meat, and as beasts of burden. Their dried droppings are an important fuel, used all over Tibet, and are often the only fuel available on the high, treeless Tibetan Plateau. Yaks transport goods across mountain passes for local farmers and traders and are an attraction for climbing and trekking expeditions: "Only one thing makes it hard to use yaks for long journeys in barren regions. They will not eat grain, which could be carried on the journey. They will starve unless they can be brought to a place where there is grass." They also are used to draw ploughs. Yaks' milk is often processed to a cheese called chhurpi in Tibetan and Nepali languages, and byaslag in Mongolia. Butter made from yaks' milk is an ingredient of the butter tea that Tibetans consume in large quantities, and is also used in lamps and made into butter sculptures used in religious festivities.
Outside the Himalayas
Small numbers of herds can be found in the United States, Canada, New Zealand, and some parts of Europe. Yaks have generated interest outside the Himalayas as a commercial crop and by cattle breeders. The main interest of North American yak breeders is lean meat production by hybridizing with other cattle, followed by wool production.
Research
The Indian government established a dedicated centre for research into yak husbandry, the ICAR-National Research Centre on Yak, in 1989. It is located at Dirang, Arunachal Pradesh, and maintains a yak farm in the Nyukmadung area at an altitude of above MSL.
Yak breeding and hybridization
In Nepal, Tibet, and Mongolia, domestic cattle are crossbred with yaks. This gives rise to the infertile male dzo མཛོ། as well as fertile females known as or མཛོ་མོ།, which may be crossed again with cattle. The Dwarf Lulu breed, "the only Bos primigenius taurus type of cattle in Nepal" has been tested for DNA markers and found to be a mixture of both taurine and zebu types of cattle (B. p. taurus and B. p. indicus) with yak. According to the International Veterinary Information Service, the low productivity of second-generation cattle–yak crosses makes them suitable only as meat animals.
Crosses between yaks and domestic cattle (Bos primigenius taurus) have been recorded in Chinese literature for at least 2,000 years. Successful crosses have also been recorded between yak and American bison, gaur, and banteng, generally with similar results to those produced with domestic cattle.
Yak domestication
Jacques et al. (2021) show that most elaborate yak-related terminologies are found within Tibetic and Gyalrongic languages. Both branches also have native terms for yak-cattle hybrids, suggesting that Tibetic and Gyalrongic speakers may have independently cross-bred yaks and cattle, predating the proto-Gyalrongic split (3221 [2169–4319] BP) from Tibeto-Gyalrongic. The oldest dated physical evidence of yak domestication is from 2,500 years BP.
Customs
In parts of Tibet and Karakorum, yak racing is a form of entertainment at traditional festivals and an important part of their culture. More recently, sports involving domesticated yaks, such as yak skiing or yak polo, are being marketed as tourist attractions in South Asian countries, including in Gilgit-Baltistan, Pakistan.
In Nepal, an annual festival is held to drink the fresh blood of yak, and it is believed that it cures various diseases such as gastritis, jaundice, and body strain. The fresh blood is extracted from the neck of a yak without killing it. The cut is healed after the ceremony is over. The ritual is believed to be originated in Tibet and Mustang.
Gallery
| Biology and health sciences | Artiodactyla | null |
230319 | https://en.wikipedia.org/wiki/Hardy%E2%80%93Weinberg%20principle | Hardy–Weinberg principle | In population genetics, the Hardy–Weinberg principle, also known as the Hardy–Weinberg equilibrium, model, theorem, or law, states that allele and genotype frequencies in a population will remain constant from generation to generation in the absence of other evolutionary influences. These influences include genetic drift, mate choice, assortative mating, natural selection, sexual selection, mutation, gene flow, meiotic drive, genetic hitchhiking, population bottleneck, founder effect, inbreeding and outbreeding depression.
In the simplest case of a single locus with two alleles denoted A and a with frequencies and , respectively, the expected genotype frequencies under random mating are for the AA homozygotes, for the aa homozygotes, and for the heterozygotes. In the absence of selection, mutation, genetic drift, or other forces, allele frequencies p and q are constant between generations, so equilibrium is reached.
The principle is named after G. H. Hardy and Wilhelm Weinberg, who first demonstrated it mathematically. Hardy's paper was focused on debunking the view that a dominant allele would automatically tend to increase in frequency (a view possibly based on a misinterpreted question at a lecture). Today, tests for Hardy–Weinberg genotype frequencies are used primarily to test for population stratification and other forms of non-random mating.
Derivation
Consider a population of monoecious diploids, where each organism produces male and female gametes at equal frequency, and has two alleles at each gene locus. We assume that the population is so large that it can be treated as infinite. Organisms reproduce by random union of gametes (the "gene pool" population model). A locus in this population has two alleles, A and a, that occur with initial frequencies and , respectively. The allele frequencies at each generation are obtained by pooling together the alleles from each genotype of the same generation according to the expected contribution from the homozygote and heterozygote genotypes, which are 1 and 1/2, respectively:
The different ways to form genotypes for the next generation can be shown in a Punnett square, where the proportion of each genotype is equal to the product of the row and column allele frequencies from the current generation.
The sum of the entries is , as the genotype frequencies must sum to one.
Note again that as , the binomial expansion of gives the same relationships.
Summing the elements of the Punnett square or the binomial expansion, we obtain the expected genotype proportions among the offspring after a single generation:
These frequencies define the Hardy–Weinberg equilibrium. It should be mentioned that the genotype frequencies after the first generation need not equal the genotype frequencies from the initial generation, e.g. . However, the genotype frequencies for all future times will equal the Hardy–Weinberg frequencies, e.g. for . This follows since the genotype frequencies of the next generation depend only on the allele frequencies of the current generation which, as calculated by equations () and (), are preserved from the initial generation:
For the more general case of dioecious diploids [organisms are either male or female] that reproduce by random mating of individuals, it is necessary to calculate the genotype frequencies from the nine possible matings between each parental genotype (AA, Aa, and aa) in either sex, weighted by the expected genotype contributions of each such mating. Equivalently, one considers the six unique diploid-diploid combinations:
and constructs a Punnett square for each, so as to calculate its contribution to the next generation's genotypes. These contributions are weighted according to the probability of each diploid-diploid combination, which follows a multinomial distribution with . For example, the probability of the mating combination is and it can only result in the genotype: . Overall, the resulting genotype frequencies are calculated as:
As before, one can show that the allele frequencies at time equal those at time , and so, are constant in time. Similarly, the genotype frequencies depend only on the allele frequencies, and so, after time are also constant in time.
If in either monoecious or dioecious organisms, either the allele or genotype proportions are initially unequal in either sex, it can be shown that constant proportions are obtained after one generation of random mating. If dioecious organisms are heterogametic and the gene locus is located on the X chromosome, it can be shown that if the allele frequencies are initially unequal in the two sexes [e.g., XX females and XY males, as in humans], in the heterogametic sex 'chases' in the homogametic sex of the previous generation, until an equilibrium is reached at the weighted average of the two initial frequencies.
Deviations from Hardy–Weinberg equilibrium
The seven assumptions underlying Hardy–Weinberg equilibrium are as follows:
organisms are diploid
only sexual reproduction occurs
generations are nonoverlapping
mating is random
population size is infinitely large
allele frequencies are equal in the sexes
there is no migration, gene flow, admixture, mutation or selection
Violations of the Hardy–Weinberg assumptions can cause deviations from expectation. How this affects the population depends on the assumptions that are violated.
Random mating. The HWP states the population will have the given genotypic frequencies (called Hardy–Weinberg proportions) after a single generation of random mating within the population. When the random mating assumption is violated, the population will not have Hardy–Weinberg proportions. A common cause of non-random mating is inbreeding, which causes an increase in homozygosity for all genes.
If a population violates one of the following four assumptions, the population may continue to have Hardy–Weinberg proportions each generation, but the allele frequencies will change over time.
Selection, in general, causes allele frequencies to change, often quite rapidly. While directional selection eventually leads to the loss of all alleles except the favored one (unless one allele is dominant, in which case recessive alleles can survive at low frequencies), some forms of selection, such as balancing selection, lead to equilibrium without loss of alleles.
Mutation will have a very subtle effect on allele frequencies through the introduction of new allele into a population. Mutation rates are of the order 10−4 to 10−8, and the change in allele frequency will be, at most, the same order. Recurrent mutation will maintain alleles in the population, even if there is strong selection against them.
Migration genetically links two or more populations together. In general, allele frequencies will become more homogeneous among the populations. Some models for migration inherently include nonrandom mating (Wahlund effect, for example). For those models, the Hardy–Weinberg proportions will normally not be valid.
Small population size can cause a random change in allele frequencies. This is due to a sampling effect, and is called genetic drift. Sampling effects are most important when the allele is present in a small number of copies.
In real world genotype data, deviations from Hardy–Weinberg Equilibrium may be a sign of genotyping error.
Sex linkage
Where the A gene is sex linked, the heterogametic sex (e.g., mammalian males; avian females) have only one copy of the gene (and are termed hemizygous), while the homogametic sex (e.g., human females) have two copies. The genotype frequencies at equilibrium are p and q for the heterogametic sex but p2, 2pq and q2 for the homogametic sex.
For example, in humans red–green colorblindness is an X-linked recessive trait. In western European males, the trait affects about 1 in 12, (q = 0.083) whereas it affects about 1 in 200 females (0.005, compared to q2 = 0.007), very close to Hardy–Weinberg proportions.
If a population is brought together with males and females with a different allele frequency in each subpopulation (males or females), the allele frequency of the male population in the next generation will follow that of the female population because each son receives its X chromosome from its mother. The population converges on equilibrium very quickly.
Generalizations
The simple derivation above can be generalized for more than two alleles and polyploidy.
Generalization for more than two alleles
Consider an extra allele frequency, r. The two-allele case is the binomial expansion of (p + q)2, and thus the three-allele case is the trinomial expansion of (p + q + r)2.
More generally, consider the alleles A1, ..., An given by the allele frequencies p1 to pn;
giving for all homozygotes:
and for all heterozygotes:
Generalization for polyploidy
The Hardy–Weinberg principle may also be generalized to polyploid systems, that is, for organisms that have more than two copies of each chromosome. Consider again only two alleles. The diploid case is the binomial expansion of:
and therefore the polyploid case is the binomial expansion of:
where c is the ploidy, for example with tetraploid (c = 4):
Whether the organism is a 'true' tetraploid or an amphidiploid will determine how long it will take for the population to reach Hardy–Weinberg equilibrium.
Complete generalization
For distinct alleles in -ploids, the genotype frequencies in the Hardy–Weinberg equilibrium are given by individual terms in the multinomial expansion of :
Significance tests for deviation
Testing deviation from the HWP is generally performed using Pearson's chi-squared test, using the observed genotype frequencies obtained from the data and the expected genotype frequencies obtained using the HWP. For systems where there are large numbers of alleles, this may result in data with many empty possible genotypes and low genotype counts, because there are often not enough individuals present in the sample to adequately represent all genotype classes. If this is the case, then the asymptotic assumption of the chi-squared distribution, will no longer hold, and it may be necessary to use a form of Fisher's exact test, which requires a computer to solve. More recently a number of MCMC methods of testing for deviations from HWP have been proposed (Guo & Thompson, 1992; Wigginton et al. 2005)
Example chi-squared test for deviation
This data is from E. B. Ford (1971) on the scarlet tiger moth, for which the phenotypes of a sample of the population were recorded. Genotype–phenotype distinction is assumed to be negligibly small. The null hypothesis is that the population is in Hardy–Weinberg proportions, and the alternative hypothesis is that the population is not in Hardy–Weinberg proportions.
From this, allele frequencies can be calculated:
and
So the Hardy–Weinberg expectation is:
Pearson's chi-squared test states:
There is 1 degree of freedom (degrees of freedom for test for Hardy–Weinberg proportions are # genotypes − # alleles). The 5% significance level for 1 degree of freedom is 3.84, and since the χ2 value is less than this, the null hypothesis that the population is in Hardy–Weinberg frequencies is not rejected.
Fisher's exact test (probability test)
Fisher's exact test can be applied to testing for Hardy–Weinberg proportions. Since the test is conditional on the allele frequencies, p and q, the problem can be viewed as testing for the proper number of heterozygotes. In this way, the hypothesis of Hardy–Weinberg proportions is rejected if the number of heterozygotes is too large or too small. The conditional probabilities for the heterozygote, given the allele frequencies are given in Emigh (1980) as
where n11, n12, n22 are the observed numbers of the three genotypes, AA, Aa, and aa, respectively, and n1 is the number of A alleles, where .
An example
Using one of the examples from Emigh (1980), we can consider the case where n = 100, and p = 0.34. The possible observed heterozygotes and their exact significance level is given in Table 4.
Using this table, one must look up the significance level of the test based on the observed number of heterozygotes. For example, if one observed 20 heterozygotes, the significance level for the test is 0.007. As is typical for Fisher's exact test for small samples, the gradation of significance levels is quite coarse.
However, a table like this has to be created for every experiment, since the tables are dependent on both n and p.
Equivalence tests
The equivalence tests are developed in order to establish sufficiently good agreement of the observed genotype frequencies and Hardy Weinberg equilibrium. Let denote the family of the genotype distributions under the assumption of Hardy Weinberg equilibrium. The distance between a genotype distribution and Hardy Weinberg equilibrium is defined by , where is some distance. The equivalence test problem is given by and , where is a tolerance parameter. If the hypothesis can be rejected then the population is close to Hardy Weinberg equilibrium with a high probability. The equivalence tests for the biallelic case are developed among others in Wellek (2004). The equivalence tests for the case of multiple alleles are proposed in Ostrovski (2020).
Inbreeding coefficient
The inbreeding coefficient, (see also F-statistics), is one minus the observed frequency of heterozygotes over that expected from Hardy–Weinberg equilibrium.
where the expected value from Hardy–Weinberg equilibrium is given by
For example, for Ford's data above:
For two alleles, the chi-squared goodness of fit test for Hardy–Weinberg proportions is equivalent to the test for inbreeding, .
The inbreeding coefficient is unstable as the expected value approaches zero, and thus not useful for rare and very common alleles. For: ; is undefined.
History
Mendelian genetics were rediscovered in 1900. However, it remained somewhat controversial for several years as it was not then known how it could cause continuous characteristics. Udny Yule (1902) argued against Mendelism because he thought that dominant alleles would increase in the population. The American William E. Castle (1903) showed that without selection, the genotype frequencies would remain stable. Karl Pearson (1903) found one equilibrium position with values of p = q = 0.5. Reginald Punnett, unable to counter Yule's point, introduced the problem to G. H. Hardy, a British mathematician, with whom he played cricket. Hardy was a pure mathematician and held applied mathematics in some contempt; his view of biologists' use of mathematics comes across in his 1908 paper where he describes this as "very simple":
To the Editor of Science: I am reluctant to intrude in a discussion concerning matters of which I have no expert knowledge, and I should have expected the very simple point which I wish to make to have been familiar to biologists. However, some remarks of Mr. Udny Yule, to which Mr. R. C. Punnett has called my attention, suggest that it may still be worth making...
Suppose that Aa is a pair of Mendelian characters, A being dominant, and that in any given generation the number of pure dominants (AA), heterozygotes (Aa), and pure recessives (aa) are as p:2q:r. Finally, suppose that the numbers are fairly large, so that mating may be regarded as random, that the sexes are evenly distributed among the three varieties, and that all are equally fertile. A little mathematics of the multiplication-table type is enough to show that in the next generation the numbers will be as (p + q)2:2(p + q)(q + r):(q + r)2, or as p1:2q1:r1, say.
The interesting question is: in what circumstances will this distribution be the same as that in the generation before? It is easy to see that the condition for this is q2 = pr. And since q12 = p1r1, whatever the values of p, q, and r may be, the distribution will in any case continue unchanged after the second generation
The principle was thus known as Hardy's law in the English-speaking world until 1943, when Curt Stern pointed out that it had first been formulated independently in 1908 by the German physician Wilhelm Weinberg. William Castle in 1903 also derived the ratios for the special case of equal allele frequencies, and it is sometimes (but rarely) called the Hardy–Weinberg–Castle Law.
Derivation of Hardy's equations
Hardy's statement begins with a recurrence relation for the frequencies p, 2q, and r. These recurrence relations follow from fundamental concepts in probability, specifically independence, and conditional probability. For example, consider the probability of an offspring from the generation being homozygous dominant. Alleles are inherited independently from each parent. A dominant allele can be inherited from a homozygous dominant parent with probability 1, or from a heterozygous parent with probability 0.5. To represent this reasoning in an equation, let represent inheritance of a dominant allele from a parent. Furthermore, let and represent potential parental genotypes in the preceding generation.
The same reasoning, applied to the other genotypes yields the two remaining recurrence relations. Equilibrium occurs when each proportion is constant between subsequent generations. More formally, a population is at equilibrium at generation when
, , and
By solving these equations necessary and sufficient conditions for equilibrium to occur can be determined. Again, consider the frequency of homozygous dominant animals. Equilibrium implies
First consider the case, where , and note that it implies that and . Now consider the remaining case, where :
where the final equality holds because the allele proportions must sum to one. In both cases, . It can be shown that the other two equilibrium conditions imply the same equation. Together, the solutions of the three equilibrium equations imply sufficiency of Hardy's condition for equilibrium. Since the condition always holds for the second generation, all succeeding generations have the same proportions.
Numerical example
Estimation of genotype distribution
An example computation of the genotype distribution given by Hardy's original equations is instructive. The phenotype distribution from Table 3 above will be used to compute Hardy's initial genotype distribution. Note that the p and q values used by Hardy are not the same as those used above.
As checks on the distribution, compute
and
For the next generation, Hardy's equations give
Again as checks on the distribution, compute
and
which are the expected values. The reader may demonstrate that subsequent use of the second-generation values for a third generation will yield identical results.
Estimation of carrier frequency
The Hardy–Weinberg principle can also be used to estimate the frequency of carriers of an autosomal recessive condition in a population based on the frequency of suffers.
Let us assume an estimated babies are born with cystic fibrosis, this is about the frequency of homozygous individuals observed in Northern European populations. We can use the Hardy–Weinberg equations to estimate the carrier frequency, the frequency of heterozygous individuals, .
As is small we can take p, , to be 1.
We therefore estimate the carrier rate to be , which is about the frequency observed in Northern European populations.
This can be simplified to the carrier frequency being about twice the square root of the birth frequency.
Graphical representation
It is possible to represent the distribution of genotype frequencies for a bi-allelic locus within a population graphically using a de Finetti diagram. This uses a triangular plot (also known as trilinear, triaxial or ternary plot) to represent the distribution of the three genotype frequencies in relation to each other. It differs from many other such plots in that the direction of one of the axes has been reversed. The curved line in the diagram is the Hardy–Weinberg parabola and represents the state where alleles are in Hardy–Weinberg equilibrium. It is possible to represent the effects of natural selection and its effect on allele frequency on such graphs. The de Finetti diagram was developed and used extensively by A. W. F. Edwards in his book Foundations of Mathematical Genetics.
| Biology and health sciences | Basics_4 | Biology |
230330 | https://en.wikipedia.org/wiki/Lungfish | Lungfish | Lungfish are freshwater vertebrates belonging to the class Dipnoi. Lungfish are best known for retaining ancestral characteristics within the Osteichthyes, including the ability to breathe air, and ancestral structures within Sarcopterygii, including the presence of lobed fins with a well-developed internal skeleton. Lungfish represent the closest living relatives of the tetrapods (which includes living amphibians, reptiles, birds and mammals). The mouths of lungfish typically bear tooth plates, which are used to crush hard shelled organisms.
Today there are only six known species of lungfish, living in Africa, South America, and Australia, though they were formerly globally distributed. The fossil record of the group extends into the Early Devonian, over 410 million years ago. The earliest known members of the group were marine, while almost all post-Carboniferous representatives inhabit freshwater environments.
Etymology
Modern Latin from the Greek δίπνοος (dipnoos) with two breathing structures, from δι- twice and πνοή breathing, breath.
Anatomy and morphology
All lungfish demonstrate an uninterrupted cartilaginous notochord and an extensively developed palatal dentition. Basal ("primitive") lungfish groups may retain marginal teeth and an ossified braincase, but derived lungfish groups, including all modern species, show a significant reduction in the marginal bones and a cartilaginous braincase. The bones of the skull roof in primitive lungfish are covered in a mineralized tissue called cosmine, but in post-Devonian lungfishes, the skull roof lies beneath the skin and the cosmine covering is lost. All modern lungfish show significant reductions and fusions of the bones of the skull roof, and the specific bones of the skull roof show no homology to the skull roof bones of ray-finned fishes or tetrapods. During the breeding season, the South American lungfish develops a pair of feathery appendages that are actually highly modified pelvic fins. These fins are thought to improve gas exchange around the fish's eggs in its nest.
Through convergent evolution, lungfishes have evolved internal nostrils similar to the tetrapods' choana, and a brain with certain similarities to Lissamphibian brain (except for the Queensland lungfish, which branched off in its own direction about 277 million years ago and has a brain resembling that of the Latimeria).
The dentition of lungfish is different from that of any other vertebrate group. "Odontodes" on the palate and lower jaws develop in a series of rows to form a fan-shaped occlusion surface. These odontodes then wear to form a uniform crushing surface. In several groups, including the modern lepidosireniformes, these ridges have been modified to form occluding blades.
The modern lungfishes have a number of larval features, which suggest paedomorphosis. They also demonstrate the largest genome among the vertebrates.
Modern lungfish all have an elongate body with fleshy, paired pectoral and pelvic fins and a single unpaired caudal fin replacing the dorsal, caudal and anal fins of most fishes.
Lungs
Lungfish have a highly specialized respiratory system. They have a distinct feature in that their lungs are connected to the larynx and pharynx without a trachea. While other species of fish can breathe air using modified, vascularized gas bladders, these bladders are usually simple sacs, devoid of complex internal structure. In contrast, the lungs of lungfish are subdivided into numerous smaller air sacs, maximizing the surface area available for gas exchange.
Most extant lungfish species have two lungs, with the exception of the Australian lungfish, which has only one. The lungs of lungfish are homologous to the lungs of tetrapods. As in tetrapods and bichirs, the lungs extend from the ventral surface of the esophagus and gut.
Perfusion of water
Of extant lungfish, only the Australian lungfish can breathe through its gills without needing air from its lung. In other species, the gills are too atrophied to allow for adequate gas exchange. When a lungfish is obtaining oxygen from its gills, its circulatory system is configured similarly to the common fish. The spiral valve of the conus arteriosus is open, the bypass arterioles of the third and fourth gill arches (which do not actually have gills) are shut, the second, fifth and sixth gill arch arterioles are open, the ductus arteriosus branching off the sixth arteriole is open, and the pulmonary arteries are closed. As the water passes through the gills, the lungfish uses a buccal pump. Flow through the mouth and gills is unidirectional. Blood flow through the secondary lamellae is countercurrent to the water, maintaining a more constant concentration gradient.
Perfusion of air
When breathing air, the spiral valve of the conus arteriosus closes (minimizing the mixing of oxygenated and deoxygenated blood), the third and fourth gill arches open, the second and fifth gill arches close (minimizing the possible loss of the oxygen obtained in the lungs through the gills), the sixth arteriole's ductus arteriosus is closed, and the pulmonary arteries open. Importantly, during air breathing, the sixth gill is still used in respiration; deoxygenated blood loses some of its carbon dioxide as it passes through the gill before reaching the lung. This is because carbon dioxide is more soluble in water. Air flow through the mouth is tidal, and through the lungs it is bidirectional and observes "uniform pool" diffusion of oxygen.
Ecology and life history
Lungfish are omnivorous, feeding on fish, insects, crustaceans, worms, mollusks, amphibians and plant matter. They have an intestinal spiral valve rather than a true stomach.
African and South American lungfish are capable of surviving seasonal drying out of their habitats by burrowing into mud and estivating throughout the dry season. Changes in physiology allow it to slow its metabolism to as little as one sixtieth of the normal metabolic rate, and protein waste is converted from ammonia to less-toxic urea (normally, lungfish excrete nitrogenous waste as ammonia directly into the water).
Burrowing is seen in at least one group of fossil lungfish, the Gnathorhizidae.
Lungfish can be extremely long-lived. A Queensland lungfish called "Granddad" at the Shedd Aquarium in Chicago was part of the permanent live collection from 1933 to 2017 after a previous residence at the Sydney Aquarium; at 109 years old, it was euthanized following a decline in health consistent with old age.
As of 2022, the oldest lungfish, and probably the oldest aquarium fish in the world is "Methuselah", an Australian lungfish long and weighing around . Methuselah is believed to be female, unlike its namesake, and is estimated to be over 90 years old.
Evolution
About 420 million years ago, during the Devonian, the last common ancestor of both lungfish and the tetrapods split into two separate evolutionary lineages, with the ancestor of the extant coelacanths diverging a little earlier from a sarcopterygian progenitor. Youngolepis and Diabolepis, dating to 419–417 million years ago, during Early Devonian (Lochkovian), are the currently oldest known lungfish, and show that the lungfishes had adapted to a diet including hard-shelled prey (durophagy) very early in their evolution. The earliest lungfish were marine. Almost all post-Carboniferous lungfish inhabit or inhabited freshwater environments. There were likely at least two transitions amongst lungfish from marine to freshwater habitats. The last common ancestor of all living lungfish likely lived sometime between the Late Carboniferous and the Jurassic. Lungfish remained present in the northern Laurasian landmasses into the Cretaceous period.
Extant lungfish
The Queensland lungfish, Neoceratodus forsteri, is endemic to Australia. Fossil records of this group date back 380 million years, around the time when the higher vertebrate classes were beginning to evolve. Fossils of lungfish belonging to the genus Neoceratodus have been uncovered in northern New South Wales, indicating that the Queensland lungfish has existed in Australia for at least 100 million years, making it a living fossil and one of the oldest living vertebrate genera on the planet. It is the most primitive surviving member of the ancient air-breathing lungfish (Dipnoi) lineages. The five other freshwater lungfish species, four in Africa and one in South America, are very different morphologically to N. forsteri. The Queensland lungfish can live for several days out of the water if it is kept moist, but will not survive total water depletion, unlike its African counterparts.
The South American lungfish, Lepidosiren paradoxa, is the single species of lungfish found in swamps and slow-moving waters of the Amazon, Paraguay, and lower Paraná River basins in South America. Notable as an obligate air-breather, it is the sole member of its family native to the Americas. Relatively little is known about the South American lungfish, or scaly salamander-fish. When immature it is spotted with gold on a black background. In the adult this fades to a brown or gray color. Its tooth-bearing premaxillary and maxillary bones are fused like other lungfish. South American lungfishes also share an autostylic jaw suspension (where the palatoquadrate is fused to the cranium) and powerful adductor jaw muscles with the extant lungfish (Dipnoi). Like the African lungfishes, this species has an elongate, almost eel-like body. It may reach a length of . The pectoral fins are thin and threadlike, while the pelvic fins are somewhat larger, and set far back. The fins are connected to the shoulder by a single bone, which is a marked difference from most fish, whose fins usually have at least four bones at their base; and a marked similarity with nearly all land-dwelling vertebrates. They have the lowest aquatic respiration of all extant lungfish species, and their gills are greatly reduced and essentially non-functional in the adults.
The marbled lungfish, Protopterus aethiopicus, is found in Africa. The marbled lungfish is smooth, elongated, and cylindrical with deeply embedded scales. The tail is very long and tapers at the end. They are the largest of the African lungfish species as they can reach a length of up to 200 cm. The pectoral and pelvic fins are also very long and thin, almost spaghetti-like. The newly hatched young have branched external gills much like those of newts. After 2 to 3 months the young transform (called metamorphosis) into the adult form, losing the external gills for gill openings. These fish have a yellowish gray or pinkish toned ground color with dark slate-gray splotches, creating a marbling or leopard effect over the body and fins. The color pattern is darker along the top and lighter below. The marbled lungfish's genome contains 133 billion base pairs, making it the largest known genome of any vertebrate. The only organisms known to have more base pairs are the protist Polychaos dubium and the flowering plant Paris japonica at 670 billion and 150 billion, respectively.
The gilled lungfish, Protopterus amphibius is a species of lungfish found in East Africa. It generally reaches only long, making it the smallest extant lungfish in the world. This lungfish is uniform blue, or slate grey in colour. It has small or inconspicuous black spots, and a pale grey belly.
The west African lungfish, Protopterus annectens, is a species of lungfish found in West Africa. It has a prominent snout and small eyes. Its body is long and eel-like, some 9–15 times the length of the head. It has two pairs of long, filamentous fins. The pectoral fins have a basal fringe and are about three times the head length, while its pelvic fins are about twice the head length. In general, three external gills are inserted posterior to the gill slits and above the pectoral fins. It has cycloid scales embedded in the skin. There are 40–50 scales between the operculum and the anus and 36–40 around the body before the origin of the dorsal fin. It has 34–37 pairs of ribs. The dorsal side is olive or brown in color and the ventral side is lighter, with great blackish or brownish spots on the body and fins except on its belly. They reach a length of about 100 cm in the wild.
The spotted lungfish, Protopterus dolloi, is a species of lungfish found in Africa. Specifically, it is found in the Kouilou-Niari Basin of the Republic of the Congo and Ogowe River basin in Gabon. It is also found in the lower and Middle Congo River Basins. Protopterus dolloi can aestivate on land by surrounding itself in a layer of dried mucus. It can reach a length of up to 130 cm.
Taxonomy
The relationship of lungfishes to the rest of the bony fish is well understood:
Lungfishes are most closely related to Powichthys, and then to the Porolepiformes.
Together, these taxa form the Dipnomorpha, the sister group to the Tetrapodomorpha.
Together, these form the Rhipidistia, the sister group to the coelacanths.
Recent molecular genetic analyses strongly support a sister relationship of lungfishes and tetrapods (Rhipidistia), with coelacanths branching slightly earlier.
The relationships among lungfishes are significantly more difficult to resolve. While Devonian lungfish had enough bone in the skull to determine relationships, post-Devonian lungfish are represented entirely by skull roofs and teeth, as the rest of the skull is cartilaginous. Additionally, many of the taxa already identified may not be monophyletic.
Phylogeny after Kemp, Cavin & Guinot, 2017
Cladogram after Brownstein et al. 2023
| Biology and health sciences | Fishes | null |
230342 | https://en.wikipedia.org/wiki/Keystone%20species | Keystone species | A keystone species is a species that has a disproportionately large effect on its natural environment relative to its abundance. The concept was introduced in 1969 by the zoologist Robert T. Paine. Keystone species play a critical role in maintaining the structure of an ecological community, affecting many other organisms in an ecosystem and helping to determine the types and numbers of various other species in the community. Without keystone species, the ecosystem would be dramatically different or cease to exist altogether. Some keystone species, such as the wolf and lion, are also apex predators.
The role that a keystone species plays in its ecosystem is analogous to the role of a keystone in an arch. While the keystone is under the least pressure of any of the stones in an arch, the arch still collapses without it. Similarly, an ecosystem may experience a dramatic shift if a keystone species is removed, even though that species was a small part of the ecosystem by measures of biomass or productivity.
It became a popular concept in conservation biology, alongside flagship and umbrella species. Although the concept is valued as a descriptor for particularly strong inter-species interactions, and has allowed easier communication between ecologists and conservation policy-makers, it has been criticized for oversimplifying complex ecological systems.
History
The concept of the keystone species was introduced in 1969 by zoologist Robert T. Paine. Paine developed the concept to explain his observations and experiments on the relationships between marine invertebrates of the intertidal zone (between the high and low tide lines), including starfish and mussels. He removed the starfish from an area, and documented the effects on the ecosystem. In his 1966 paper, Food Web Complexity and Species Diversity, Paine had described such a system in Makah Bay in Washington.
In his 1969 paper, Paine proposed the keystone species concept, using Pisaster ochraceus, a species of starfish generally known as ochre starfish, and Mytilus californianus, a species of mussel, as a primary example. The ochre starfish is a generalist predator and feeds on chitons, limpets, snails, barnacles, echinoids, and even decapod crustacea. The favourite food for these starfish is the mussel which is a dominant competitor for the space on the rocks. The ochre starfish keeps the population numbers of the mussels in check along with the other preys allowing the other seaweeds, sponges, and anemones, that ochre starfish do not consume, to co-exist. When Paine removed the ochre starfish, the mussels quickly outgrew the other species crowding them out. At the start, the rock pools held 15 rock-clinging species. Three years later there were 8 such species; and ten years later the pools were largely occupied by a single species, mussels. The concept became popular in conservation, and was deployed in a range of contexts and mobilized to engender support for conservation, especially where human activities had damaged ecosystems, such as by removing keystone predators.
Definitions
A keystone species was defined by Paine as a species that has a disproportionately large effect on its environment relative to its abundance. It has been defined operationally by Davic in 2003 as "a strongly interacting species whose top-down effect on species diversity and competition is large relative to its biomass dominance within a functional group."
A classic keystone species is a predator that prevents a particular herbivorous species from eliminating dominant plant species. If prey numbers are low, keystone predators can be even less abundant and still be effective. Yet without the predators, the herbivorous prey would explode in numbers, wipe out the dominant plants, and dramatically alter the character of the ecosystem. The exact scenario changes in each example, but the central idea remains that through a chain of interactions, a non-abundant species has an outsized impact on ecosystem functions. For example, the herbivorous weevil Euhrychiopsis lecontei is thought to have keystone effects on aquatic plant diversity by foraging on nuisance Eurasian watermilfoil in North American waters. Similarly, the wasp species Agelaia vicina has been labeled a keystone species for its unparalleled nest size, colony size, and high rate of brood production. The diversity of its prey and the quantity necessary to sustain its high rate of growth have a direct impact on other species around it.
The keystone concept is defined by its ecological effects, and these in turn make it important for conservation. In this it overlaps with several other species conservation concepts such as flagship species, indicator species, and umbrella species. For example, the jaguar is a charismatic big cat which meets all of these definitions:
Predators
Sea otters and kelp forests
Sea otters protect kelp forests from damage by sea urchins. When the sea otters of the North American west coast were hunted commercially for their fur, their numbers fell to such low levels – fewer than 1000 in the north Pacific ocean – that they were unable to control the sea urchin population. The urchins, in turn, grazed the holdfasts of kelp so heavily that the kelp forests largely disappeared, along with all the species that depended on them. Reintroducing the sea otters has enabled the kelp ecosystem to be restored. For example, in Southeast Alaska some 400 sea otters were released, and they have bred to form a population approaching 25,000.
The wolf, Yellowstone's apex predator
Keystone predators may increase the biodiversity of communities by preventing a single species from becoming dominant. They can have a profound influence on the balance of organisms in a particular ecosystem. Introduction or removal of a keystone predator, or changes in its population density, can have drastic cascading effects on the equilibrium of many other populations in the ecosystem. For example, grazers of a grassland may prevent a single dominant species from taking over.
The elimination of the gray wolf from the Greater Yellowstone Ecosystem had profound impacts on the trophic pyramid. Without predation, herbivores began to over-graze many woody browse species, affecting the area's plant populations. In addition, wolves often kept animals from grazing in riparian areas, which protected beavers from having their food sources encroached upon. The removal of wolves had a direct effect on beaver populations, as their habitat became grazing territory. Increased browsing on willows and conifers along Blacktail Creek due to a lack of predation caused channel incision because the beavers helped slow the water down, allowing soil to stay in place. Furthermore, predation keeps hydrological features such as creeks and streams in normal working order. When wolves were reintroduced, the beaver population and the whole riparian ecosystem recovered dramatically within a few years.
Sea stars and other non-apex predators
As described by Paine in 1966, some sea stars (e.g., Pisaster ochraceus) may prey on sea urchins, mussels, and other shellfish that have no other natural predators. If the sea star is removed from the ecosystem, the mussel population explodes uncontrollably, driving out most other species.
These creatures need not be apex predators. Sea stars are prey for sharks, rays, and sea anemones. Sea otters are prey for orca.
The jaguar, whose numbers in Central and South America have been classified as near threatened, acts as a keystone predator by its widely varied diet, helping to balance the mammalian jungle ecosystem with its consumption of 87 different species of prey. The lion is another keystone species.
Mutualists
Keystone mutualists are organisms that participate in mutually beneficial interaction, the loss of which would have a profound impact upon the ecosystem as a whole. For example, in the Avon Wheatbelt region of Western Australia, there is a period of each year when Banksia prionotes (acorn banksia) is the sole source of nectar for honeyeaters, which play an important role in pollination of numerous plant species. Therefore, the loss of this one species of tree would probably cause the honeyeater population to collapse, with profound implications for the entire ecosystem. Another example is frugivores, such as the cassowary, which spreads the seeds of many different trees. Some seeds will not grow unless they have been through a cassowary.
Ecosystem engineers
A term used alongside keystone is ecosystem engineer. In North America, the prairie dog is an ecosystem engineer. Prairie dog burrows provide the nesting areas for mountain plovers and burrowing owls. Prairie dog tunnel systems also help channel rainwater into the water table to prevent runoff and erosion, and can also serve to change the composition of the soil in a region by increasing aeration and reversing soil compaction that can be a result of cattle grazing. Prairie dogs also trim the vegetation around their colonies, perhaps to remove any cover for predators. Grazing species such as plains bison, which is another keystone species, the pronghorn, and the mule deer have shown a proclivity for grazing on the same land used by prairie dogs.
The beaver is a well known ecosystem engineer and keystone species. It transforms its territory from a stream to a pond or swamp. Beavers affect the environment first altering the edges of riparian areas by cutting down older trees to use for their dams. This allows younger trees to take their place. Beaver dams alter the riparian area they are established in. Depending on topography, soils, and many factors, these dams change the riparian edges of streams and rivers into wetlands, meadows, or riverine forests. These dams have been shown to be beneficial to a myriad of species including amphibians, salmon, and song birds.
In the African savanna, the larger herbivores, especially the elephants, shape their environment. The elephants destroy trees, making room for the grass species and creating habitat for various small animal species. Without these animals, much of the savanna would turn into woodland.
In the Amazon river basin, peccaries produce and maintain wallows that are utilized by a wide variety of species.
Australian studies have found that parrotfish on the Great Barrier Reef are the only reef fish that consistently scrape and clean the coral on the reef. Without these animals, the Great Barrier Reef would be under severe strain.
In the Serengeti, the presence of sufficient gnus in these grasslands reduces wildfire likelihood, which in turn promotes tree growth. The documentary The Serengeti Rules documents this in detail.
Limitations
Depends on context
The community ecologist Bruce Menge states that the keystone concept has been stretched far beyond Paine's original concept. That stretching can be quantified: the researcher Ishana Shukla has listed 230 species identified as keystones in some 157 studies in the 50 years since Paine's paper. Menge's own work has shown that the purple Pisaster sea star that Paine had studied was a powerful keystone species in places exposed to strong wave action, but was far less important in sheltered places. Paine had indeed stated that in Alaska, without the relevant mussel species as prey, the predatory Pisaster was "just another sea star". In other words, the extent to which a species could be described as a keystone depended on the ecological context.
Multiple meanings
Although the concept of the keystone species has a value in describing particularly strong inter-species interactions, and for allowing easier communication between ecologists and conservation policy-makers, it has been criticized by L. S. Mills and colleagues for oversimplifying complex ecological systems. The term has been applied widely in different ecosystems and to predators, prey, and plants (primary producers), inevitably with differing ecological meanings. For instance, removing a predator may allow other animals to increase to the point where they wipe out other species; removing a prey species may cause predator populations to crash, or may allow predators to drive other prey species to extinction; and removing a plant species may result in the loss of animals that depend on it, like pollinators and seed dispersers. Beavers too have been called keystone, not for eating other species but for modifying the environment in ways that affected other species. The term has thus been given quite different meanings in different cases. In Mills's view, Paine's work showed that a few species could sometimes have extremely strong interactions within a particular ecosystem, but that does not automatically imply that other ecosystems have a similar structure.
| Biology and health sciences | Ecology | Biology |
230361 | https://en.wikipedia.org/wiki/Fin%20whale | Fin whale | The fin whale (Balaenoptera physalus), also known as the finback whale or common rorqual, is a species of baleen whale and the second-longest cetacean after the blue whale. The biggest individual reportedly measured in length, with a maximum recorded weight of . The fin whale's body is long, slender and brownish-gray in color, with a paler underside to appear less conspicuous from below (countershading).
At least two recognized subspecies exist, one in the North Atlantic and one across the Southern Hemisphere. It is found in all the major oceans, from polar to tropical waters, though it is absent only from waters close to the pack ice at the poles and relatively small areas of water away from the open ocean. The highest population density occurs in temperate and cool waters. Its prey mainly consists of smaller schooling fish, small squid, or crustaceans, including copepods and krill. Mating takes place in temperate, low-latitude seas during the winter. Fin whales are often observed in pods of 6–10 animals, with whom they communicate utilizing frequency-modulated sounds, ranging from 16 to 40 hertz.
Like all other large whales, the fin whale was a prized kill during the "heyday" of whaling, from 1840 to 1861. It remained so into the 20th century but decades of over harvesting contributed to declining numbers through the late 20th century. Over 725,000 fin whales were reportedly taken from the Southern Hemisphere between 1905 and 1976. Post-recovery numbers of the southern subspecies are predicted to be less than 50% of the pre-whaling population, even by 2100, due to long-lasting impacts of whaling and slow recovery rates. As of 2018, it was assessed as vulnerable by the IUCN.
Taxonomy
The fin whale was first described by Friderich Martens in 1675 and by Paul Dudley in 1725. The former description was used as the primary basis for the species Balaena physalus by Carl Linnaeus in 1758. In 1804, Bernard Germain de Lacépède reclassified the species as Balaenoptera rorqual, based on a specimen that had stranded on Île Sainte-Marguerite (Cannes, France) in 1798. In 1830, Louis Companyo described a specimen that had been stranded near Saint-Cyprien, southern France, in 1828 as Balaena musculus. Most later authors followed him in using the specific name musculus, until Frederick W. True (1898) showed that it referred to the blue whale. In 1846, British taxonomist John Edward Gray described a specimen from the Falkland Islands as Balaenoptera australis. In 1865, German naturalist Hermann Burmeister described a roughly specimen found near Buenos Aires about 30 years earlier as Balaenoptera patachonicus. In 1903, Romanian scientist Emil Racoviță placed all these designations into Balaenoptera physalus. The word physalus comes from the Greek word physa, meaning "blows," referring to the prominent blow of the species.
Fin whales are rorquals, members of the family Balaenopteridae, which includes the humpback whale, the blue whale, Bryde's whale, the sei whale, and the minke whale. The family diverged from the other baleen whales in the suborder Mysticeti as long ago as the middle Miocene.
Recent DNA evidence indicates the fin whale may be more closely related to the humpback whale (Megaptera novaeangliae) and, in at least one study, the gray whale (Eschrichtius robustus), two whales in different genera, than it is to members of its own genus, such as the minke whales.
As of 2023, four subspecies are named, each with distinct physical features and vocalizations. The northern fin whale, B. p. physalus (Linnaeus 1758) inhabits the North Atlantic and the southern fin whale, B. p. quoyi (Fischer 1829) occupies the Southern Hemisphere. Most experts consider the fin whales of the North Pacific to be a third subspecies—this was supported by a 2013 study, which found that the Northern Hemisphere B. p. physalus was not composed of a single subspecies. A 2019 genetic study concluded that the North Pacific fin whales should be considered a subspecies, suggesting the name B. p. velifera (Scammon 1869). The three groups mix at most rarely.
Clarke (2004) proposed a "pygmy" subspecies (B. p. patachonica, Burmeister, 1865) that is purportedly darker in colour and has black baleen. He based this on a single physically mature female caught in the Antarctic in 1947–48, the smaller average size (a few feet) of sexually and physically mature fin whales caught by the Japanese around 50°S, and smaller, darker sexually immature fin whales caught in the Antarctic which he believed were a "migratory phase" of his proposed subspecies. The subspecies has not been genetically established, and is not recognized by the Society for Marine Mammalogy.
Hybrids
The genetic distance between blue and fin whales has been compared to that between a chimpanzee and human (3.5 million years on the evolutionary tree.) Nevertheless, hybrid individuals between blue and fin whales with characteristics of both are known to occur with relative frequency in both the North Atlantic and North Pacific.
The DNA profile of a sampling of whale meat in the Japanese market found evidence of blue/fin hybrids. Similarly, a whale caught by whalers off the coast of Iceland in 2018 was found to be a hybrid descended from a female blue whale and a male fin whale. A 2024 genome analysis of North Atlantic blue whales found that approximately 3.5% of their genome was derived from hybridization with fin whales. The gene flow was determined to be unidirectional from fin to blue whales. Despite their smaller size, fin whales have similar cruising and sprinting speeds to blue whales, which would allow fin males to complete courtship chases with blue females.
Anatomy
The body is relatively thin with a slender rostrum and large hook-like dorsal fin that is situated in the posterior fourth of the body. It has an elongated ridge on its back, and around 350 to 400 baleen plates. Like all rorquals, the fin whale has grooves between the tip of the lower jaw and the navel.
Among whale species, the fin whale is exceeded in size only by the blue whale. Adults usually average 40 to 50 tonnes in weight. Males have a mean length of , and females of . They are sexually dimorphic, with females generally being longer and heavier than males. The largest specimens can attain lengths of over and weights of 77 to 81 tonnes.
The fin whale is brownish to dark or light gray dorsally and white ventrally. The left side of the head is dark gray, while the right side exhibits a complex pattern of contrasting light and dark markings. The right lower jaw is white or light gray, which sometimes extends laterally and dorsally unto the upper jaw. Dark, oval-shaped areas of pigment called "flipper shadows" extend below and posterior to the pectoral fins.
The penis size of fin whales typically reaches a length of ; the testes usually weigh in mature individuals. The oral cavity of the fin whale has a very stretchy or extensible nerve system which aids them in feeding.
Life history
Mating takes place during the winter months, in temperate, low-latitude waters, and the gestation period lasts between 11 and 12 months. At 6 or 7 months of age, when it is in length, a newborn weans from its mother, and the calf accompanies its mother to the summer feeding area. Although reports of up to six foetuses have been made, single births are far more typical. Females reproduce every two to three years. In the Northern Hemisphere, females reach sexual maturity between the ages of 6 and 12 at lengths of , and around in the Southern Hemisphere. Calves remain with their mothers for about a year.
Full physical maturity is attained between 25 and 30 years. Fin whales have a maximum life span of at least 94 years of age, although specimens have been found aged at an estimated 135–140 years. The fin whale is one of the fastest cetaceans and can sustain speeds between and and bursts up to have been recorded, earning the fin whale the nickname "the greyhound of the sea". Fin whales are more gregarious than other rorquals, and often live in groups of 6–10, although feeding groups may reach up to 100 animals.
Vocalizations
Like other whales, males make long, loud, low-frequency sounds. The vocalizations of blue and fin whales are the lowest-frequency sounds made by any animal. Most sounds are frequency-modulated (FM) down-swept infrasonic pulses from 16 to 40 hertz frequency (the range of sounds that most humans can hear falls between 20 hertz and 20 kilohertz). Each sound lasts one to two seconds, and various sound combinations occur in patterned sequences lasting 7 to 15 minutes each. The whale then repeats the sequences in bouts lasting up to many days. The vocal sequences have source levels of up to 184–186 decibels relative to 1 micropascal at a reference distance of one metre and can be detected hundreds of miles from their source.
When fin whale sounds were first recorded by US biologists, they did not realize that these unusually loud, long, pure and regular sounds were being made by whales. They first investigated the possibilities that the sounds were due to equipment malfunction, geophysical phenomena, or even part of a Soviet Union scheme for detecting enemy submarines. Eventually, biologists demonstrated that the sounds were the vocalizations of fin whales.
Direct association of these vocalizations with the reproductive season for the species and that only males make the sounds point to these vocalizations as possible reproductive displays. Over the past 100 years, the dramatic increase in ocean noise from shipping and naval activity may have slowed the recovery of the fin whale population, by impeding communications between males and receptive females. Fin whale songs can penetrate over below the sea floor and seismologists can use those song waves to assist in underwater surveys.
Breathing
When feeding, fin whales blow five to seven times in quick succession, but while traveling or resting will blow once every minute or two. On their terminal (last) dive they arch their back high out of the water, but rarely raise their flukes out of the water. They then dive to depths of up to when feeding or a few hundred feet when resting or traveling. The average feeding dive off California and Baja lasts 6 minutes, with a maximum of 17 minutes; when traveling or resting they usually dive for only a few minutes at a time.
Ecology
Range and habitat
Like many large rorquals, the fin whale is a cosmopolitan species. It is found in all the world's major oceans and in waters ranging from the polar to the tropical. It is absent only from waters close to the ice pack at both the north and south extremities and relatively small areas of water away from the large oceans, such as the Red Sea, although they can reach into the Baltic Sea, a marginal sea of such conditions. The highest population density occurs in temperate and cool waters. It is less densely populated in the warmest, equatorial regions.
The North Atlantic fin whale has an extensive distribution, occurring from the Gulf of Mexico and Mediterranean Sea, northward to Baffin Bay and Spitsbergen. In general, fin whales are more common north of approximately 30°N latitude, but considerable confusion arises about their occurrence south of 30°N latitude because of the difficulty in distinguishing fin whales from Bryde's whales. Extensive ship surveys have led researchers to conclude that the summer feeding range of fin whales in the western North Atlantic is mainly between 41°20'N and 51°00'N, from shore seaward to the contour.
Summer distribution of fin whales in the North Pacific is the immediate offshore waters from central Baja California to Japan and as far north as the Chukchi Sea bordering the Arctic Ocean. They occur in high densities in the northern Gulf of Alaska and southeastern Bering Sea between May and October, with some movement through the Aleutian passes into and out of the Bering Sea. Several whales tagged between November and January off southern California were killed in the summer off central California, Oregon, British Columbia, and in the Gulf of Alaska. Fin whales have been observed feeding 250 miles south of Hawaii in mid-May, and several winter sightings have been made there. Some researchers have suggested that the whales migrate into Hawaiian waters primarily in the autumn and winter.
Although fin whales are certainly migratory, moving seasonally in and out of high-latitude feeding areas, the overall migration pattern is not well understood. Acoustic readings from passive-listening hydrophone arrays indicate a southward migration of the North Atlantic fin whale occurs in the autumn from the Labrador-Newfoundland region, south past Bermuda, and into the West Indies. One or more populations of fin whales are thought to remain year-round in high latitudes, moving offshore, but not southward in late autumn. A study based on resightings of identified fin whales in Massachusetts Bay indicates that calves often learn migratory routes from their mothers and return to their mother's feeding area in subsequent years.
In the Pacific, migration patterns are poorly characterized. Although some fin whales are apparently present year-round in the Gulf of California, there is a significant increase in their numbers in the winter and spring. Southern fin whales migrate seasonally from relatively high-latitude Antarctic feeding grounds in the summer to low-latitude breeding and calving areas in the winter. The location of winter breeding areas is still unknown, since these whales tend to migrate in the open ocean.
It has been shown that populations of fin whales within the Mediterranean have preferred feeding locations that partially overlap with high concentrations of plastic pollution and microplastic debris. High concentrations of microplastics most likely overlap with fin whales' preferred feeding grounds because both microplastic and the whale's food sources are near high trophic upwelling areas.
The total historical North Pacific population was estimated at 42,000 to 45,000 before the start of whaling. Of this, the population in the eastern portion of the North Pacific was estimated to be 25,000 to 27,000. Surveys conducted in 1991, 1993, 1996, and 2001 produced estimates between 1,600 and 3,200 off California and 280 and 380 off Oregon and Washington. Surveys in coastal waters of British Columbia in summers 2004 and 2005 produced abundance estimates of approximately 500 animals. Fin whales might have started returning to the coastal waters off British Columbia (a sighting occurred in Johnstone Strait in 2011) and Kodiak Island. Size of the local population migrating to Hawaiian Archipelago is unknown.
Finbacks are also relatively abundant along the coast of Peru and Chile (in Chile, most notably off Los Lagos region such as Gulf of Corcovado in Chiloé National Park, , port of Mejillones, and Caleta Zorra. Year-round confirmations indicate possible residents off pelagic north eastern to central Chile such as around coastal Caleta Chañaral and Pingüino de Humboldt National Reserve, east of Juan Fernández Islands, and northeast of Easter Island and possible wintering ground exist for eastern south Pacific population.
Among Northern Indian Ocean and Bay of Bengal, such as along Sri Lanka, India, and Malaysia, sightings and older records of fin whales exist.
Predation
The only known predator of the fin whale is the killer whale, with at least 20 eyewitness and second-hand accounts of attack or harassment. They usually flee and offer little resistance to attack. Only a few confirmed fatalities have occurred. In October 2005, 16 killer whales attacked and killed a fin whale in the Canal de Ballenas, Gulf of California, after chasing it for about an hour. They fed on its sinking carcass for about 15 minutes before leaving the area. In June 2012, a pod of killer whales was seen in La Paz Bay, in the Gulf of California, chasing a fin whale for over an hour before finally killing it and feeding on its carcass. The whale bore numerous tooth rakes over its back and dorsal fin; several killer whales flanked it on either side, with one individual visible under water biting at its right lower jaw. In July 1908, a whaler reportedly saw two killer whales attack and kill a fin whale off western Greenland. In January 1984, seven were seen from the air circling, holding the flippers, and ramming a fin whale in the Gulf of California, but the observation ended at nightfall.
Feeding
The fin whale is a filter-feeder, feeding on small schooling fish, squid and crustaceans including copepods and krill. In the North Pacific, they feed on krill in the genera Euphausia, Thysanoessa, and Nyctiphanes, large copepods in the genus Neocalanus, small schooling fish (e.g. the genera Engraulis, Mallotus, Clupea, and Theragra), and squid. Based on stomach content analysis of over 19,500 fin whales caught by the Japanese whaling fleet in the North Pacific from 1952 to 1971, 64.1% contained only krill, 25.5% copepods, 5.0% fish, 3.4% krill and copepods and 1.7% squid. Nemoto (1959) analyzed the stomach contents of about 7500 fin whales caught in the northern North Pacific and Bering Sea from 1952 to 1958, found that they mainly preyed on euphausiids around the Aleutian Islands and in the Gulf of Alaska and schooling fish in the northern Bering Sea and off Kamchatka.
Of the fin whale stomachs sampled off British Columbia between 1963 and 1967, euphausiids dominated the diet for four of the five years (82.3 to 100% of the diet), while copepods only formed a major portion of the diet in 1965 (35.7%). Miscellaneous fish, squid, and octopus played only a very minor part of the diet in two of the five years (3.6 to 4.8%). Fin whales caught off California between 1959 and 1970 fed on the pelagic euphausiid Euphausia pacifica (86% of sampled individuals), the more neritic euphausiid Thysanoessa spinifera (9%), and the northern anchovy (Engraulis mordax) (7%); only trace amounts (<0.5% each) were found of Pacific saury (C. saira) and juvenile rockfish (Sebastes jordani).
In the North Atlantic, they prey on euphausiids in the genera Meganyctiphanes, Thysanoessa and Nyctiphanes and small schooling fish (e.g. the genera Clupea, Mallotus, and Ammodytes). Of the 1,609 fin whale stomachs examined at the Hvalfjörður whaling station in southwestern Iceland from 1967 to 1989 (caught between June and September), 96% contained only krill, 2.5% krill and fish, 0.8% some fish remains, 0.7% capelin (M. villosus), and 0.1% sandeel (family Ammodytidae); a small proportion of (mainly juvenile) blue whiting (Micromesistius poutassou) were also found. Of the krill sampled between 1979 and 1989, the vast majority (over 99%) was northern krill (Meganyctiphanes norvegica); only one stomach contained Thysanoessa longicaudata. Off West Greenland, 75% of the fin whales caught between July and October had consumed krill (family Euphausiidae), 17% capelin (Mallotus) and 8% sand lance (Ammodytes sp.). Off eastern Newfoundland, they chiefly feed on capelin, but also take small quantities of euphausiids (mostly T. raschii and T. inermis). In the Ligurian-Corsican-Provençal Basin in the Mediterranean Sea they make dives as deep as to feed on the euphausiid Meganyctiphanes norvegica, while off the island of Lampedusa, between Tunisia and Sicily, they have been observed in mid-winter feeding on surface swarms of the small euphausiid Nyctiphanes couchi.
In the Southern Hemisphere, they feed almost exclusively on euphausiids (mainly the genera Euphausia and Thysanoessa), as well as taking small amounts of amphipods (e.g. Themisto gaudichaudii) and various species of fish. Of the more than 16,000 fin whales caught by the Japanese whaling fleet in the Southern Hemisphere between 1961 and 1965 that contained food in their stomachs, 99.4% fed on euphausiids, 0.5% on fish, and 0.1% on amphipods. In the Southern Ocean they mainly consume E. superba.
The animal feeds by opening its jaws while swimming at some in one study, which causes it to engulf up to of water in one gulp. It then closes its jaws and pushes the water back out of its mouth through its baleen, which allows the water to leave while trapping the prey. An adult has between 262 and 473 baleen plates on each side of the mouth. Each plate is made of keratin that frays out into fine hairs on the ends inside the mouth near the tongue. Each plate can measure up to in length and in width.
The whale routinely dives to depths of more than where it executes an average of four "lunges", to accumulate krill. Each gulp provides the whale with approximately of food. One whale can consume up to of food a day, leading scientists to conclude that the whale spends about three hours a day feeding to meet its energy requirements, roughly the same as humans. If prey patches are not sufficiently dense, or are located too deep in the water, the whale has to spend a larger portion of its day searching for food. One hunting technique is to circle schools of fish at high speed, frightening the fish into a tight ball, then turning on its side before engulfing the massed prey.
Parasites, epibiotics, and pathology
Fin whales suffer from a number of pathological conditions. The parasitic copepod Pennella balaenopterae—usually found on the flank of fin whales—burrows into their blubber to feed on their blood, while the pseudo-stalked barnacle Xenobalanus globicipitis is generally found more often on the dorsal fin, pectoral fins, and flukes.
Other barnacles found on fin whales include the acorn barnacle Coronula reginae and the stalked barnacle Conchoderma auritum, which attaches to Coronula or the baleen. The harpacticid copepod Balaenophilus unisetus (heavy infestations of which have been found in fin whales caught off northwestern Spain) and the ciliate Haematophagus also infest the baleen, the former feeding on the baleen itself and the latter on red blood cells.
The remora Remora australis and occasionally the amphipod Cyamus balaenopterae can also be found on fin whales, both feeding on the skin. Infestations of the giant nematode Crassicauda boopis can cause inflammation of the renal arteries and potential kidney failure, while the smaller C. crassicauda infects the lower urinary tract. Out of 87 whales taken and necropsied from the North Atlantic, infection from Crassicauda boopis was found to be very prevalent and invasive, indicating high probability that it was responsible for causing death in these whales. C. boopis was found in 94% of the whales examined. The worms were usually enveloped by "exuberant tissue reactions which in some whales obstructed multiple renal veins". The parasite was most likely by environmental contamination, involving shedding of larvae in urine. Major inflammatory lesions in the mesenteric arteries suggested that the worm larvae were ingested and migrated to the kidney.
These observations suggest that infection from C. boopis can be "lethal by inducing congestive renal failure". Injury to the vascular system is also a result of moderate infections. Therefore, the implication can be made that the feeding migration of fin whales every year in circumpolar waters can be associated with pathologic risk.
An emaciated female fin whale, which stranded along the Belgian coast in 1997, was found to be infected with lesions of Morbillivirus. In January 2011, a emaciated adult male fin whale stranded dead on the Tyrrhenian coastline of Italy was found to be infected with Morbillivirus and the protozoa Toxoplasma gondii, as well as carrying heavy loads of organochlorine pollutants.
Human interaction
Whaling
In the 19th century, the fin whale was occasionally hunted by open-boat whalers, but it was relatively safe, because it could easily outrun ships of the time and often sank when killed, making the pursuit a waste of time for whalers. However, the later introduction of steam-powered boats and harpoons that exploded on impact made it possible to kill and secure them along with blue and sei whales on an industrial scale. As other whale species became overhunted, the whaling industry turned to the still-abundant fin whale as a substitute. It was primarily hunted for its blubber, oil, and baleen. Around 704,000 fin whales were caught in Antarctic whaling operations alone between 1904 and 1975.
The introduction of factory ships with stern slipways in 1925 substantially increased the number of whales taken per year. By 1962–63, sei whale catches began to increase as fin whales became scarce. Coastal groups in northeast Asian waters, along with many other baleen species, were likely driven into serious perils or functional extinctions by industrial catches by Japan covering wide ranges of China and Korean EEZ within a very short period in the 20th century. After the cease of exploiting Asian stocks, Japan kept mass commercial and illegal hunts until 1975. Several thousand individuals were hunted from various stations mainly along coasts of Hokkaido, Sanriku, and the Gotō Islands.
The IWC prohibited hunting in the Southern Hemisphere in 1976. The Soviet Union engaged in the illegal killing of protected whale species in the North Pacific and Southern Hemisphere, over-reporting fin whale catches to cover up illegal takes of other species. The fin whale was given full protection from commercial whaling by the IWC in the North Pacific in 1976, and in the North Atlantic in 1987, with small exceptions for aboriginal catches and catches for research purposes. All populations worldwide remain listed as endangered species by the US National Marine Fisheries Service and the International Conservation Union Red List.
The IWC has set a quota of 19 fin whales per year for Greenland. Meat and other products from whales killed in these hunts are widely marketed within Greenland, but export is illegal. Iceland and Norway are not bound by the IWC's moratorium on commercial whaling because both countries filed objections to it.
In the Southern Hemisphere, Japan permitted annual takes of 10 fin whales under its Antarctic Special Permit whaling program for the 2005–2006 and 2006–2007 seasons. The proposal for 2007–2008 and the subsequent 12 seasons allowed takes of 50 per year. In 2019, Japan left the International Whaling Commission (IWC) and resumed commercial whaling. Japan reported a catch of 212 total whales in both 2020 and 2021; however, no fin whale catches have yet been reported.
Ship interaction
Collisions with ships are a major cause of mortality. In some areas, they cause a substantial portion of large whale strandings. Most serious injuries are caused by large, fast-moving ships over or near continental shelves.
A 60-foot-long fin whale was found stuck on the bow of a container ship in New York harbour on 12 April 2014. Two dead fin whales, one 65 feet and one 25 feet, were discovered stuck to the Australian destroyer HMAS Sydney in May 2021 when the ship arrived in Naval Base San Diego.
Ship collisions frequently occur in Tsushima Strait and result in damage done to whales, passengers, and vessels. In response the Japanese Coast Guard has started a surveillance program to monitor large cetacean activity in Tsushima Strait to inform operating vessels in the area.
Whale watching
Fin whales are regularly encountered on whale-watching excursions worldwide. In Monterey Bay and the Southern California Bight, fin whales are encountered year-round, with the best sightings between November and March. They can even be seen from land (for example, from Point Vicente, Palos Verdes, where they can be seen lunge feeding at the surface only a half mile to a few miles offshore). They are regularly sighted in the summer and fall in the Gulf of St. Lawrence, the Gulf of Maine, the Bay of Fundy, the Bay of Biscay, Strait of Gibraltar, and the Mediterranean. In southern Ireland, they are seen inshore from June to February, with peak sightings in November and December.
Conservation
As of 2018, the global fin whale population is estimated to be 100,000 mature individuals. There are an estimated total of 70,000 individuals in the North Atlantic, 50,000 in the North Pacific, and 25,000 in the Southern Hemisphere.
The fin whale is listed as a vulnerable species on the IUCN Red List. They are also included in the Endangered Species Act of 1973. The fin whale is listed on both Appendix I and Appendix II of the Convention on the Conservation of Migratory Species of Wild Animals (CMS). Commercial whaling of the species was officially banned in 1976, both in the North Pacific and Southern Hemisphere. Post-whaling populations have steadily increased. The fin whale is still hunted off the waters of West Greenland, and in the Antarctic Ocean by Japanese researchers.
They may also become entangled in fishing gear in some rare instances. Military sonar may effect the behavioral patterns of fin whales, which can lead to population decline. Similarly, whale watching may cause fin whales to alter their behavior and foraging habits.
The fin whale is covered by the Agreement on the Conservation of Cetaceans in the Black Sea, Mediterranean Sea and Contiguous Atlantic Area (ACCOBAMS) and the Memorandum of Understanding for the Conservation of Cetaceans and Their Habitats in the Pacific Islands Region (Pacific Cetaceans MOU).
| Biology and health sciences | Baleen whales | Animals |
230364 | https://en.wikipedia.org/wiki/Pumice | Pumice | Pumice (), called pumicite in its powdered or dust form, is a volcanic rock that consists of extremely vesicular rough-textured volcanic glass, which may or may not contain crystals. It is typically light-colored. Scoria is another vesicular volcanic rock that differs from pumice in having larger vesicles, thicker vesicle walls, and being dark colored and denser.
Pumice is created when super-heated, highly pressurized rock is rapidly ejected from a volcano. The unusual foamy configuration of pumice happens because of simultaneous rapid cooling and rapid depressurization. The depressurization creates bubbles by lowering the solubility of gases (including water and CO2) that are dissolved in the lava, causing the gases to rapidly exsolve (like the bubbles of CO2 that appear when a carbonated drink is opened). The simultaneous cooling and depressurization freeze the bubbles in a matrix. Eruptions under water are rapidly cooled and the large volume of pumice created can be a shipping hazard for cargo ships.
Properties
Pumice is composed of highly microvesicular glass pyroclastic with very thin, translucent bubble walls of extrusive igneous rock. It is commonly but not exclusively of silicic or felsic to intermediate in composition (e.g., rhyolitic, dacitic, andesite, pantellerite, phonolite, trachyte), but basaltic and other compositions are known. Pumice is commonly pale in color, ranging from white, cream, blue or grey, to green-brown or black. It forms when volcanic gases exsolving from viscous magma form bubbles that remain within the viscous magma as it cools to glass. Pumice is a common product of explosive eruptions (plinian and ignimbrite-forming) and commonly forms zones in upper parts of silicic lavas. Pumice has a porosity of 64–85% by volume and it floats on water, possibly for years, until it eventually becomes waterlogged and sinks.
Scoria differs from pumice in being denser. With larger vesicles and thicker vesicle walls, scoria sinks rapidly. The difference is the result of the lower viscosity of the magma that forms scoria. When larger amounts of gas are present, the result is a finer-grained variety of pumice known as pumicite. Pumicite consists of particles less than in size. Pumice is considered a volcanic glass because it has no crystal structure. Pumice varies in density according to the thickness of the solid material between the bubbles; many samples float in water. After the explosion of Krakatoa, rafts of pumice drifted through the Indian Ocean for up to 20 years, with tree trunks floating among them. In fact, pumice rafts disperse and support several marine species. In 1979, 1984 and 2006, underwater volcanic eruptions near Tonga created large pumice rafts that floated hundreds of kilometres to Fiji.
There are two main forms of vesicles. Most pumice contains tubular microvesicles that can impart a silky or fibrous fabric. The elongation of the microvesicles occurs due to ductile elongation in the volcanic conduit or, in the case of pumiceous lavas, during flow. The other form of vesicles are subspherical to spherical and result from high vapor pressure during an eruption. Reticulite is a type of basaltic pumice formed in very high lava fountains. It has an extremely low density and is composed of a network of volcanic glass formed when the vesicles have almost completely coalesced.
Etymology
Pumice is an igneous rock with a foamy appearance. The name is derived from the Latin word pumex (meaning "pumice") which is related to the Latin word spuma meaning "foam". In former times, pumice was called "Spuma Maris", meaning "froth of the sea" in Latin because the frothy material was thought to be hardened sea foam. Around 80 B.C., it was called "lapis spongiae" in Latin for its vesicular properties. Many Greek scholars decided there were different sources of pumice, one of which was in the sea coral category.
Area
Pumice can be found all around the globe deriving from continental volcanic occurrences and submarine volcanic occurrences. Floating stones can also be distributed by ocean currents. As described earlier pumice is produced by the eruption of explosive volcanoes under certain conditions, therefore, natural sources occur in volcanically active regions. Pumice is mined and transported from these regions. In 2011, Italy and Turkey led pumice mining production at 4 and 3 million tons respectively; other large producers at or exceeding a million tonnes were Greece, Iran, Chile, and Syria. Total world pumice production in 2011 was estimated at 17 million tonnes.
Asia
There are large reserves of pumice in Asian countries including Afghanistan, Indonesia, Japan, Syria, Iran, and eastern Russia. Considerable amounts of pumice can be found at the Kamchatka Peninsula on the eastern flank of Russia. This area contains 19 active volcanoes and it lies in close proximity with the Pacific volcanic belt. Asia is also the site of the second-most dangerous volcanic eruption in the 20th century, Mount Pinatubo, which erupted on June 12, 1991 in the Philippines. Ash and pumice lapilli were distributed over a mile around the volcano. These ejections filled trenches that once reached 660 feet deep. So much magma was displaced from the vent that the volcano became a depression on the surface of the Earth. Another well-known volcano that produces pumice is Krakatoa. An eruption in 1883 ejected so much pumice that kilometers of sea were covered in floating pumice and in some areas rose 1.5 meters above sea level.
Europe
Europe is the largest producer of pumice with deposits in Italy, Turkey, Greece, Hungary, Iceland, and Germany. Italy is the largest producer of pumice because of its numerous eruptive volcanoes. On the Aeolian Islands of Italy, the island of Lipari is entirely made up of volcanic rock, including pumice. Large amounts of igneous rock on Lipari are due to the numerous extended periods of volcanic activity from the Late Pleistocene (Tyrrhenian) to the Holocene.
North America
Pumice can be found all across North America including on the Caribbean Islands. In the United States, pumice is mined in Nevada, Oregon, Idaho, Arizona, California, New Mexico and Kansas. U.S. production of pumice and pumicite in 2011 was estimated at 380,000 tonnes, valued at $7.7 million with approximately 46% coming from Nevada and Oregon. Idaho is also known as a large producer of pumice because of the quality and brightness of the rock found in local reserves. One of the most famous volcanoes was Mount Mazama that erupted 7,700 years ago in Oregon and deposited 300 feet of pumice and ash around the vent. The large amount of magma that was erupted caused the structure to collapse, forming a caldera now known as Crater Lake.
South America
Chile is one of the leading producers of pumice in the world. The Puyehue-Cordón Caulle are two coalesced volcanoes in the Andes mountains that ejected ash and pumice across Chile and Argentina. A recent eruption in 2011 wreaked havoc on the region by covering all surfaces and lakes in ash and pumice.
Africa
Kenya, Ethiopia and Tanzania have some deposits of pumice.
New Zealand
The Havre Seamount volcano produced the largest-known deep ocean volcanic eruption on Earth. The volcano erupted in July 2012 but remained unnoticed until enormous pieces of pumice were seen to be floating on the Pacific Ocean. Blankets of rock reached a thickness of 5 meters. Most of this floating pumice is deposited on the northwest coast of New Zealand and the Polynesia islands.
Mining
The mining of pumice is an environmentally friendly process compared with other mining methods because the igneous rock is deposited on the surface of the earth in loose aggregate form. The material is mined by open-pit methods. Soils are removed by machinery in order to obtain more pure quality pumice. Scalping screens are used to filter impure surficial pumice of organic soils and unwanted rocks. Blasting is not necessary because the material is unconsolidated, therefore only simple machinery is used such as bulldozers and power shovels. Different sizes of pumice are needed for specific uses therefore crushers are used to achieve desired grades ranging from lump, coarse, intermediate, fine, and extra fine.
Uses
Pumice is a very lightweight, porous and abrasive material and it has been used for centuries in the construction and beauty industry as well as in early medicine. It is also used as an abrasive, especially in polishes, pencil erasers, and the production of stone-washed jeans. Pumice was also used in the early book-making industry to prepare parchment paper and leather bindings. There is high demand for pumice, particularly for water filtration, chemical spill containment, cement manufacturing, horticulture and increasingly for the pet industry. The mining of pumice in environmentally sensitive areas has been under more scrutiny after such an operation was stopped in the U.S. state of Oregon, at Rock Mesa in the southern part of the Three Sisters Wilderness.
Early medicine
Pumice has been used in the medicinal industry for more than 2000 years. Ancient Chinese medicine used ground pumice along with ground mica and fossilized bones added to teas to calm the spirit. This tea was used to treat dizziness, nausea, insomnia, and anxiety disorders. Ingestion of these pulverized rocks was believed to be able to soften nodules and was later used with other herbal ingredients to treat gallbladder cancer and urinary difficulties. In Western medicine, beginning in the early 18th century, pumice ground into a sugar consistency mixed with other ingredients was used to attempt to treat ulcers mostly on the skin and cornea. Concoctions such as these were also used to help wounds scar in a supposedly healthier manner. In approximately 1680 it was noted by an English naturalist that pumice powder was used to promote sneezing.
Personal care
Pumice has been used as a material in personal care for thousands of years. It is an abrasive material that can be used in powdered form or as a stone to remove unwanted hair or skin. In ancient Egypt, it was common to remove all hair on the body to control lice and as a form of ritual purification, using creams, razors, and pumice stones. Pumice in powdered form was an ingredient in toothpastes in ancient Rome. Nail care was very important in ancient China; nails were kept groomed with pumice stones, and pumice stones were also used to remove calluses.
It was discovered in a Roman poem that pumice was used to remove dead skin as far back as 100 BC, and likely before then. It has been used throughout many eras since then, including the Victorian Era. Today, many of these techniques are still used; pumice is widely used as a skin exfoliant. Even though hair removal techniques have evolved over the centuries, abrasive material like pumice stones is also still used. "Pumice stones" are often used in beauty salons during the pedicure process to remove dry and excess skin from the bottom of the foot as well as calluses.
Finely ground pumice has been added to some toothpastes as a polish, similar to Roman use, and easily removes dental plaque build-up. Such toothpaste is too abrasive for daily use. Pumice is also added to heavy-duty hand cleaners (such as lava soap) as a mild abrasive. Some brands of chinchilla dust bath are formulated with powdered pumice. Old beauty techniques using pumice are still employed today but newer substitutes are easier to obtain.
Cleaning
Pumice stone, sometimes attached to a handle, is an effective scrubbing tool for removal of limescale, rust, hard water rings, and other stains on porcelain fixtures in households (e.g., bathrooms). It is a quick method compared to alternatives like chemicals or vinegar and baking soda or borax.
Horticulture
Good soil requires sufficient water and nutrient loading as well as little compaction to allow easy exchange of gases. The roots of plants require continuous transportation of carbon dioxide and oxygen to and from the surface. Pumice improves the quality of soil because of its porous properties; water and gases can be transported easily through the pores and nutrients can be stored in the microscopic holes. Pumice rock fragments are inorganic therefore no decomposition and little compaction occur.
Another benefit of this inorganic rock is that it does not attract or host fungi or insects. As drainage is very important in horticulture, with the presence of pumice tillage is much easier. Pumice usage also creates ideal conditions for growing plants like cacti and succulents as it increases the water retention in sandy soils and reduces the density of clayey soils to allow more transportation of gases and water. The addition of pumice to soil improves and increases vegetative cover as the roots of plants make slopes more stable therefore it helps reduce erosion. It is often used on roadsides and ditches and commonly used in turf and golf courses to maintain grass cover and flatness that can degrade due to large amounts of traffic and compaction. Chemically pumice is pH neutral, neither acidic nor alkaline. In 2011, 16% of pumice mined in the United States was used for horticultural purposes.
Pumice contributes to soil fertility in areas where it is naturally present in the soil due to volcanic activity. For example, in the Jemez Mountains of New Mexico, the Ancestral Puebloans settled on "pumice patches" of the El Cajete Pumice which likely retained a greater amount of moisture and was ideal for farming.
Construction
Pumice is widely used to make lightweight concrete and insulative low-density cinder blocks. The air-filled vesicles in this porous rock serve as a good insulator. A fine-grained version of pumice called pozzolan is used as an additive in cement and is mixed with lime to form a light-weight, smooth, plaster-like concrete. This form of concrete was used as far back as Roman times. Roman engineers utilized it to build the huge dome of the Pantheon with increasing amounts of pumice added to concrete for higher elevations of the structure. It was also commonly used as a construction material for many aqueducts.
One of the main uses of pumice currently in the United States is manufacturing concrete. This rock has been used in concrete mixtures for thousands of years and continues to be used in producing concrete, especially in regions close to where this volcanic material is deposited.
New studies prove a broader application of pumice powder in the concrete industry. Pumice can act as a cementitious material in concrete and researchers have shown that concrete made with up to 50% pumice powder can significantly improve durability yet reduce greenhouse gas emissions and fossil fuel consumption.
| Physical sciences | Igneous rocks | Earth science |
230401 | https://en.wikipedia.org/wiki/Floyd%E2%80%93Warshall%20algorithm | Floyd–Warshall algorithm | In computer science, the Floyd–Warshall algorithm (also known as Floyd's algorithm, the Roy–Warshall algorithm, the Roy–Floyd algorithm, or the WFI algorithm) is an algorithm for finding shortest paths in a directed weighted graph with positive or negative edge weights (but with no negative cycles). A single execution of the algorithm will find the lengths (summed weights) of shortest paths between all pairs of vertices. Although it does not return details of the paths themselves, it is possible to reconstruct the paths with simple modifications to the algorithm. Versions of the algorithm can also be used for finding the transitive closure of a relation , or (in connection with the Schulze voting system) widest paths between all pairs of vertices in a weighted graph.
History and naming
The Floyd–Warshall algorithm is an example of dynamic programming, and was published in its currently recognized form by Robert Floyd in 1962. However, it is essentially the same as algorithms previously published by Bernard Roy in 1959 and also by Stephen Warshall in 1962 for finding the transitive closure of a graph, and is closely related to Kleene's algorithm (published in 1956) for converting a deterministic finite automaton into a regular expression, with the difference being the use of a min-plus semiring. The modern formulation of the algorithm as three nested for-loops was first described by Peter Ingerman, also in 1962.
Algorithm
The Floyd–Warshall algorithm compares many possible paths through the graph between each pair of vertices. It is guaranteed to find all shortest paths and is able to do this with comparisons in a graph, even though there may be edges in the graph. It does so by incrementally improving an estimate on the shortest path between two vertices, until the estimate is optimal.
Consider a graph with vertices numbered 1 through . Further consider a function that returns the length of the shortest possible path (if one exists) from to using vertices only from the set as intermediate points along the way. Now, given this function, our goal is to find the length of the shortest path from each to each using any vertex in . By definition, this is the value , which we will find recursively.
Observe that must be less than or equal to : we have more flexibility if we are allowed to use the vertex . If is in fact less than , then there must be a path from to using the vertices that is shorter than any such path that does not use the vertex . Since there are no negative cycles this path can be decomposed as:
(1) a path from to that uses the vertices , followed by
(2) a path from to that uses the vertices .
And of course, these must be a shortest such path (or several of them), otherwise we could further decrease the length. In other words, we have arrived at the recursive formula:
.
The base case is given by
where denotes the weight of the edge from to if one exists and ∞ (infinity) otherwise.
These formulas are the heart of the Floyd–Warshall algorithm. The algorithm works by first computing for all pairs for , then , then , and so on. This process continues until , and we have found the shortest path for all pairs using any intermediate vertices. Pseudocode for this basic version follows.
Pseudocode
let dist be a |V| × |V| array of minimum distances initialized to ∞ (infinity)
for each edge (u, v) do
dist[u][v] = w(u, v) // The weight of the edge (u, v)
for each vertex v do
dist[v][v] = 0
for k from 1 to |V|
for i from 1 to |V|
for j from 1 to |V|
if dist[i][j] > dist[i][k] + dist[k][j]
dist[i][j] = dist[i][k] + dist[k][j]
end if
Example
The algorithm above is executed on the graph on the left below:
Prior to the first recursion of the outer loop, labeled above, the only known paths correspond to the single edges in the graph. At , paths that go through the vertex 1 are found: in particular, the path [2,1,3] is found, replacing the path [2,3] which has fewer edges but is longer (in terms of weight). At , paths going through the vertices {1,2} are found. The red and blue boxes show how the path [4,2,1,3] is assembled from the two known paths [4,2] and [2,1,3] encountered in previous iterations, with 2 in the intersection. The path [4,2,3] is not considered, because [2,1,3] is the shortest path encountered so far from 2 to 3. At , paths going through the vertices {1,2,3} are found. Finally, at , all shortest paths are found.
The distance matrix at each iteration of , with the updated distances in bold, will be:
Behavior with negative cycles
A negative cycle is a cycle whose edges sum to a negative value. There is no shortest path between any pair of vertices , which form part of a negative cycle, because path-lengths from to can be arbitrarily small (negative). For numerically meaningful output, the Floyd–Warshall algorithm assumes that there are no negative cycles. Nevertheless, if there are negative cycles, the Floyd–Warshall algorithm can be used to detect them. The intuition is as follows:
The Floyd–Warshall algorithm iteratively revises path lengths between all pairs of vertices , including where ;
Initially, the length of the path is zero;
A path can only improve upon this if it has length less than zero, i.e. denotes a negative cycle;
Thus, after the algorithm, will be negative if there exists a negative-length path from back to .
Hence, to detect negative cycles using the Floyd–Warshall algorithm, one can inspect the diagonal of the path matrix, and the presence of a negative number indicates that the graph contains at least one negative cycle. During the execution of the algorithm, if there is a negative cycle, exponentially large numbers can appear, as large as , where is the largest absolute value of a negative edge in the graph. To avoid overflow/underflow problems one should check for negative numbers on the diagonal of the path matrix within the inner for loop of the algorithm. Obviously, in an undirected graph a negative edge creates a negative cycle (i.e., a closed walk) involving its incident vertices. Considering all edges of the above example graph as undirected, e.g. the vertex sequence 4 – 2 – 4 is a cycle with weight sum −2.
Path reconstruction
The Floyd–Warshall algorithm typically only provides the lengths of the paths between all pairs of vertices. With simple modifications, it is possible to create a method to reconstruct the actual path between any two endpoint vertices. While one may be inclined to store the actual path from each vertex to each other vertex, this is not necessary, and in fact, is very costly in terms of memory. Instead, we can use the shortest-path tree, which can be calculated for each node in time using memory, and allows us to efficiently reconstruct a directed path between any two connected vertices.
Pseudocode
The array holds the penultimate vertex on the path from to (except in the case of , where it always contains even if there is no self-loop on ):
let dist be a array of minimum distances initialized to (infinity)
let prev be a array of vertex indices initialized to null
procedure FloydWarshallWithPathReconstruction() is
for each edge (u, v) do
dist[u][v] = w(u, v) // The weight of the edge (u, v)
prev[u][v] = u
for each vertex v do
dist[v][v] = 0
prev[v][v] = v
for k from 1 to |V| do // standard Floyd-Warshall implementation
for i from 1 to |V|
for j from 1 to |V|
if dist[i][j] > dist[i][k] + dist[k][j] then
dist[i][j] = dist[i][k] + dist[k][j]
prev[i][j] = prev[k][j]
procedure Path(u, v) is
if prev[u][v] = null then
return []
path = [v]
while u ≠ v do
v = prev[u][v]
path.prepend(v)
return path
Time complexity
Let be , the number of vertices. To find all of
(for all and ) from those of
requires operations. Since we begin with
and compute the sequence of matrices , , , , each having a cost of ,
the total time complexity of the algorithm is .
Applications and generalizations
The Floyd–Warshall algorithm can be used to solve the following problems, among others:
Shortest paths in directed graphs (Floyd's algorithm).
Transitive closure of directed graphs (Warshall's algorithm). In Warshall's original formulation of the algorithm, the graph is unweighted and represented by a Boolean adjacency matrix. Then the addition operation is replaced by logical conjunction (AND) and the minimum operation by logical disjunction (OR).
Finding a regular expression denoting the regular language accepted by a finite automaton (Kleene's algorithm, a closely related generalization of the Floyd–Warshall algorithm)
Inversion of real matrices (Gauss–Jordan algorithm)
Optimal routing. In this application one is interested in finding the path with the maximum flow between two vertices. This means that, rather than taking minima as in the pseudocode above, one instead takes maxima. The edge weights represent fixed constraints on flow. Path weights represent bottlenecks; so the addition operation above is replaced by the minimum operation.
Fast computation of Pathfinder networks.
Widest paths/Maximum bandwidth paths
Computing canonical form of difference bound matrices (DBMs)
Computing the similarity between graphs
Transitive closure in AND/OR/threshold graphs.
Implementations
Implementations are available for many programming languages.
For C++, in the boost::graph library
For C#, at QuikGraph
For C#, at QuickGraphPCL (A fork of QuickGraph with better compatibility with projects using Portable Class Libraries.)
For Java, in the Apache Commons Graph library
For JavaScript, in the Cytoscape library
For Julia, in the Graphs.jl package
For MATLAB, in the Matlab_bgl package
For Perl, in the Graph module
For Python, in the SciPy library (module scipy.sparse.csgraph) or NetworkX library
For R, in packages e1071 and Rfast
For C, a pthreads, parallelized, implementation including a SQLite interface to the data at floydWarshall.h
Comparison with other shortest path algorithms
For graphs with non-negative edge weights, Dijkstra's algorithm can be used to find all shortest paths from a single vertex with running time . Thus, running Dijkstra starting at each vertex takes time . Since , this yields a worst-case running time of repeated Dijkstra of . While this matches the asymptotic worst-case running time of the Floyd-Warshall algorithm, the constants involved matter quite a lot. When a graph is dense (i.e., ), the Floyd-Warshall algorithm tends to perform better in practice. When the graph is sparse (i.e., is significantly smaller than ), Dijkstra tends to dominate.
For sparse graphs with negative edges but no negative cycles, Johnson's algorithm can be used, with the same asymptotic running time as the repeated Dijkstra approach.
There are also known algorithms using fast matrix multiplication to speed up all-pairs shortest path computation in dense graphs, but these typically make extra assumptions on the edge weights (such as requiring them to be small integers). In addition, because of the high constant factors in their running time, they would only provide a speedup over the Floyd–Warshall algorithm for very large graphs.
| Mathematics | Graph theory | null |
230426 | https://en.wikipedia.org/wiki/BL%20Lacertae%20object | BL Lacertae object | A BL Lacertae object or BL Lac object is a type of active galactic nucleus (AGN) or a galaxy with such an AGN, named after its prototype, BL Lacertae. In contrast to other types of active galactic nuclei, BL Lacs are characterized by rapid and large-amplitude flux variability and significant optical polarization. Because of these properties, the prototype of the class (BL Lac) was originally thought to be a variable star. When compared to the more luminous active nuclei (quasars) with strong emission lines, BL Lac objects have spectra dominated by a relatively featureless non-thermal emission continuum over the entire electromagnetic range. This lack of spectral lines historically hindered identification of the nature and distance of such objects.
In the unified scheme of radio-loud active galactic nuclei, the observed nuclear phenomenology of BL Lacs is interpreted as being due to the effects of the relativistic jet closely aligned to the line of sight of the observer. BL Lacs are thought to be intrinsically identical to low-power radio galaxies. These active nuclei appear to be hosted in massive elliptical galaxies. From the point of AGN classification, BL Lacs are a blazar subtype. All known BL Lacs are associated with core dominated radio sources, many of them exhibiting apparent superluminal motion.
The blazar category encompasses all quasars oriented with the relativistic jet directed at the observer giving a unique radio emission spectrum. This includes BL Lacs as well as optically violent variable (OVV) quasars, however in general practice, "Blazar" and "BL Lac Object" are often used interchangeably. OVV quasars are generally more luminous and have stronger emission lines than BL Lac objects.
Some examples of BL Lac objects are BL Lacertae itself, OJ 287, AP Librae, PKS 2155-304, PKS 0521-365, Markarian 421, 3C 371, W Comae Berenices, ON 325 and Markarian 501.
Host galaxies
Soon after the discovery of this unusual class of objects it was noted that the sources were surrounded by a faint nebulosity. In the late 1970s the use of modern detectors (such as CCD) allowed observers to probe with better accuracy the nature of the nebulosity. First images of the BL Lac object PKS 0548-322 by Michael John Disney in 1974 in various filters found it to be composed by a giant elliptical galaxy with a bright nucleus.
Extensive surveys taken with the Hubble Space Telescope of 132 BL Lac objects comprising seven complete radio, X-ray, and optically selected samples in 2000 studied the morphologies of possible BL Lac host galaxies. The data concluded that in two-thirds of the BLL images taken, host galaxies are detected, including in nearly all with redshift z < 0.5. BL Lac objects are luminous enough that only one quarter (6/22) of the images taken with z > 0.5 were resolved because of relatively short exposure times. A de Vaucouleurs profile looks to be a significantly preferred brightness profile for 58 of the 72 resolved host galaxies at over ~99% confidence. The results of this survey conclude that there is an 8% limit to the number of disk systems in BL Lac objects and is therefore consistent with the assumption that all BL Lac host galaxies could be elliptical. These ellipticals are very luminous with a median absolute K-corrected magnitude of mag (rms dispersion). This is comparable to the brightest cluster galaxies.
History
John L. Schmitt first noticed the peculiar nature of BL Lac in 1968 when he matched it with a radio object, VRO 42.22.01.
Within a year others observed that the radio flux varied, and that light was polarized. Peter Albert Strittmatter proposed the class of object in 1972 and added four objects. By 1976 there were 30 known objects.
In 2017, a very high energy neutrino was detected by the IceCube project apparently coming from BL Lac object TXS 0506+056.
| Physical sciences | Active galactic nucleus | Astronomy |
230428 | https://en.wikipedia.org/wiki/Angular%20resolution | Angular resolution | Angular resolution describes the ability of any image-forming device such as an optical or radio telescope, a microscope, a camera, or an eye, to distinguish small details of an object, thereby making it a major determinant of image resolution. It is used in optics applied to light waves, in antenna theory applied to radio waves, and in acoustics applied to sound waves. The colloquial use of the term "resolution" sometimes causes confusion; when an optical system is said to have a high resolution or high angular resolution, it means that the perceived distance, or actual angular distance, between resolved neighboring objects is small. The value that quantifies this property, θ, which is given by the Rayleigh criterion, is low for a system with a high resolution. The closely related term spatial resolution refers to the precision of a measurement with respect to space, which is directly connected to angular resolution in imaging instruments. The Rayleigh criterion shows that the minimum angular spread that can be resolved by an image-forming system is limited by diffraction to the ratio of the wavelength of the waves to the aperture width. For this reason, high-resolution imaging systems such as astronomical telescopes, long distance telephoto camera lenses and radio telescopes have large apertures.
Definition of terms
Resolving power is the ability of an imaging device to separate (i.e., to see as distinct) points of an object that are located at a small angular distance or it is the power of an optical instrument to separate far away objects, that are close together, into individual images. The term resolution or minimum resolvable distance is the minimum distance between distinguishable objects in an image, although the term is loosely used by many users of microscopes and telescopes to describe resolving power. As explained below, diffraction-limited resolution is defined by the Rayleigh criterion as the angular separation of two point sources when the maximum of each source lies in the first minimum of the diffraction pattern (Airy disk) of the other. In scientific analysis, in general, the term "resolution" is used to describe the precision with which any instrument measures and records (in an image or spectrum) any variable in the specimen or sample under study.
The Rayleigh criterion
The imaging system's resolution can be limited either by aberration or by diffraction causing blurring of the image. These two phenomena have different origins and are unrelated. Aberrations can be explained by geometrical optics and can in principle be solved by increasing the optical quality of the system. On the other hand, diffraction comes from the wave nature of light and is determined by the finite aperture of the optical elements. The lens' circular aperture is analogous to a two-dimensional version of the single-slit experiment. Light passing through the lens interferes with itself creating a ring-shape diffraction pattern, known as the Airy pattern, if the wavefront of the transmitted light is taken to be spherical or plane over the exit aperture.
The interplay between diffraction and aberration can be characterised by the point spread function (PSF). The narrower the aperture of a lens the more likely the PSF is dominated by diffraction. In that case, the angular resolution of an optical system can be estimated (from the diameter of the aperture and the wavelength of the light) by the Rayleigh criterion defined by Lord Rayleigh: two point sources are regarded as just resolved when the principal diffraction maximum (center) of the Airy disk of one image coincides with the first minimum of the Airy disk of the other, as shown in the accompanying photos. (In the bottom photo on the right that shows the Rayleigh criterion limit, the central maximum of one point source might look as though it lies outside the first minimum of the other, but examination with a ruler verifies that the two do intersect.) If the distance is greater, the two points are well resolved and if it is smaller, they are regarded as not resolved. Rayleigh defended this criterion on sources of equal strength.
Considering diffraction through a circular aperture, this translates into:
where θ is the angular resolution (radians), λ is the wavelength of light, and D is the diameter of the lens' aperture. The factor 1.22 is derived from a calculation of the position of the first dark circular ring surrounding the central Airy disc of the diffraction pattern. This number is more precisely 1.21966989... (), the first zero of the order-one Bessel function of the first kind divided by π.
The formal Rayleigh criterion is close to the empirical resolution limit found earlier by the English astronomer W. R. Dawes, who tested human observers on close binary stars of equal brightness. The result, θ = 4.56/D, with D in inches and θ in arcseconds, is slightly narrower than calculated with the Rayleigh criterion. A calculation using Airy discs as point spread function shows that at Dawes' limit there is a 5% dip between the two maxima, whereas at Rayleigh's criterion there is a 26.3% dip. Modern image processing techniques including deconvolution of the point spread function allow resolution of binaries with even less angular separation.
Using a small-angle approximation, the angular resolution may be converted into a spatial resolution, Δℓ, by multiplication of the angle (in radians) with the distance to the object. For a microscope, that distance is close to the focal length f of the objective. For this case, the Rayleigh criterion reads:
.
This is the radius, in the imaging plane, of the smallest spot to which a collimated beam of light can be focused, which also corresponds to the size of the smallest object that the lens can resolve. The size is proportional to wavelength, λ, and thus, for example, blue light can be focused to a smaller spot than red light. If the lens is focusing a beam of light with a finite extent (e.g., a laser beam), the value of D corresponds to the diameter of the light beam, not the lens. Since the spatial resolution is inversely proportional to D, this leads to the slightly surprising result that a wide beam of light may be focused on a smaller spot than a narrow one. This result is related to the Fourier properties of a lens.
A similar result holds for a small sensor imaging a subject at infinity: The angular resolution can be converted to a spatial resolution on the sensor by using f as the distance to the image sensor; this relates the spatial resolution of the image to the f-number, #:
.
Since this is the radius of the Airy disk, the resolution is better estimated by the diameter,
Specific cases
Single telescope
Point-like sources separated by an angle smaller than the angular resolution cannot be resolved. A single optical telescope may have an angular resolution less than one arcsecond, but astronomical seeing and other atmospheric effects make attaining this very hard.
The angular resolution R of a telescope can usually be approximated by
where λ is the wavelength of the observed radiation, and D is the diameter of the telescope's objective. The resulting R is in radians. For example, in the case of yellow light with a wavelength of 580 nm, for a resolution of 0.1 arc second, we need D=1.2 m. Sources larger than the angular resolution are called extended sources or diffuse sources, and smaller sources are called point sources.
This formula, for light with a wavelength of about 562 nm, is also called the Dawes' limit.
Telescope array
The highest angular resolutions for telescopes can be achieved by arrays of telescopes called astronomical interferometers: These instruments can achieve angular resolutions of 0.001 arcsecond at optical wavelengths, and much higher resolutions at x-ray wavelengths. In order to perform aperture synthesis imaging, a large number of telescopes are required laid out in a 2-dimensional arrangement with a dimensional precision better than a fraction (0.25x) of the required image resolution.
The angular resolution R of an interferometer array can usually be approximated by
where λ is the wavelength of the observed radiation, and B is the length of the maximum physical separation of the telescopes in the array, called the baseline. The resulting R is in radians. Sources larger than the angular resolution are called extended sources or diffuse sources, and smaller sources are called point sources.
For example, in order to form an image in yellow light with a wavelength of 580 nm, for a resolution of 1 milli-arcsecond, we need telescopes laid out in an array that is 120 m × 120 m with a dimensional precision better than 145 nm.
Microscope
The resolution R (here measured as a distance, not to be confused with the angular resolution of a previous subsection) depends on the angular aperture :
where .
Here NA is the numerical aperture, is half the included angle of the lens, which depends on the diameter of the lens and its focal length, is the refractive index of the medium between the lens and the specimen, and is the wavelength of light illuminating or emanating from (in the case of fluorescence microscopy) the sample.
It follows that the NAs of both the objective and the condenser should be as high as possible for maximum resolution. In the case that both NAs are the same, the equation may be reduced to:
The practical limit for is about 70°. In a dry objective or condenser, this gives a maximum NA of 0.95. In a high-resolution oil immersion lens, the maximum NA is typically 1.45, when using immersion oil with a refractive index of 1.52. Due to these limitations, the resolution limit of a light microscope using visible light is about 200 nm. Given that the shortest wavelength of visible light is violet (),
which is near 200 nm.
Oil immersion objectives can have practical difficulties due to their shallow depth of field and extremely short working distance, which calls for the use of very thin (0.17 mm) cover slips, or, in an inverted microscope, thin glass-bottomed Petri dishes.
However, resolution below this theoretical limit can be achieved using super-resolution microscopy. These include optical near-fields (Near-field scanning optical microscope) or a diffraction technique called 4Pi STED microscopy. Objects as small as 30 nm have been resolved with both techniques. In addition to this Photoactivated localization microscopy can resolve structures of that size, but is also able to give information in z-direction (3D).
List of telescopes and arrays by angular resolution
| Physical sciences | Basics | Astronomy |
230456 | https://en.wikipedia.org/wiki/Mallard | Mallard | The mallard () or wild duck (Anas platyrhynchos) is a dabbling duck that breeds throughout the temperate and subtropical Americas, Eurasia, and North Africa. It has been introduced to New Zealand, Australia, Peru, Brazil, Uruguay, Argentina, Chile, Colombia, the Falkland Islands, and South Africa. Belonging to the subfamily Anatinae of the waterfowl family Anatidae, mallards live in wetlands, eat water plants and small animals, and are social animals preferring to congregate in groups or flocks of varying sizes.
Males (drakes) have green heads, while the females (hens) have mainly brown-speckled plumage. Both sexes have an area of white-bordered black or iridescent purple or blue feathers called a speculum on their wings; males especially tend to have blue speculum feathers. The mallard is long, of which the body makes up around two-thirds the length. The wingspan is and the bill is long. It is often slightly heavier than most other dabbling ducks, weighing .
The female lays 8 to 13 creamy white to greenish-buff spotless eggs, on alternate days. Incubation takes 27 to 28 days and fledging takes 50 to 60 days. The ducklings are precocial and fully capable of swimming as soon as they hatch.
The non-migratory mallard interbreeds with indigenous wild ducks of closely related species through genetic pollution by producing fertile offspring. Complete hybridisation of various species of wild duck gene pools could result in the extinction of many indigenous waterfowl. This species is the main ancestor of most breeds of domestic duck, and its naturally evolved wild gene pool has been genetically polluted by the domestic and feral mallard populations.
The mallard is considered to be a species of least concern by the International Union for Conservation of Nature (IUCN), and, unlike many waterfowl, are considered an invasive species in some regions. It is a very adaptable species, being able to live and even thrive in urban areas which may have supported more localised, sensitive species of waterfowl before development.
Taxonomy and evolutionary history
The mallard was one of the many bird species originally described in the 1758 10thedition of Systema Naturae by Carl Linnaeus. He gave it two binomial names: Anas platyrhynchos and Anas boschas. The latter was generally preferred until 1906 when Einar Lönnberg established that A.platyrhynchos had priority, as it appeared on an earlier page in the text. The scientific name comes from Latin Anas, "duck" and Ancient Greek πλατυρυγχος, platyrhynchus, "broad-billed" (from πλατύς, platys, "broad" and ρυγχός, rhunkhos, "bill"). The genome of Anas platyrhynchos was sequenced in 2013.
The name mallard originally referred to any wild drake, and it is sometimes still used this way. It was derived from the Old French or for "wild drake" although its true derivation is unclear. It may be related to, or at least influenced by, an Old High German masculine proper name , clues lying in the alternative English forms "maudelard" and "mawdelard". Masle (male) has also been proposed as an influence.
Mallards frequently interbreed with their closest relatives in the genus Anas, such as the American black duck, and also with species more distantly related, such as the northern pintail, leading to various hybrids that may be fully fertile. The mallard has hybridised with more than 40 species in the wild, and an additional 20 species in captivity, though fertile hybrids typically have two Anas parents. Mallards and their domestic conspecifics are fully interfertile; many wild mallard populations in North America contain significant amounts of domestic mallard DNA.
Genetic analysis has shown that certain mallards appear to be closer to their Indo-Pacific relatives, while others are related to their American relatives. Mitochondrial DNA data for the D-loop sequence suggest that mallards may have evolved in the general area of Siberia. Mallard bones rather abruptly appear in food remains of ancient humans and other deposits of fossil bones in Europe, without a good candidate for a local predecessor species. The large Ice Age palaeosubspecies that made up at least the European and West Asian populations during the Pleistocene has been named Anas platyrhynchos palaeoboschas.Mallards are differentiated in their mitochondrial DNA between North American and Eurasian populations, but the nuclear genome displays a notable lack of genetic structure. Haplotypes typical of American mallard relatives and eastern spot-billed ducks can be found in mallards around the Bering Sea. The Aleutian Islands hold a population of mallards that appear to be evolving towards becoming a subspecies, as gene flow with other populations is very limited.
Also, the paucity of morphological differences between the Old World mallards and the New World mallard demonstrates the extent to which the genome is shared among them such that birds like the Chinese spot-billed duck are highly similar to the Old World mallard, and birds such as the Hawaiian duck are highly similar to the New World mallard.
The size of the mallard varies clinally; for example, birds from Greenland, though larger, have smaller bills, paler plumage, and stockier bodies than birds further south and are sometimes classified as a separate subspecies, the Greenland mallard (A.p.conboschas).
Description
The mallard is a medium-sized waterfowl species that is often slightly heavier than most other dabbling ducks. It is longof which the body makes up around two-thirdshas a wingspan of , and weighs . Among standard measurements, the wing chord is , the bill is , and the tarsus is .The breeding male mallard is unmistakable, with a glossy bottle-green head and a white collar that demarcates the head and neck from the purple-tinged brown breast, grey-brown wings, and a pale grey belly. The rear of the male is black, with white-bordered dark tail feathers. The bill of the male is a yellowish-orange tipped with black, with that of the female generally darker and ranging from black to mottled orange and brown. The female mallard is predominantly mottled, with each individual feather showing sharp contrast from buff to very dark brown, a coloration shared by most female dabbling ducks, and has buff cheeks, eyebrow, throat, and neck, with a darker crown and eye-stripe. Mallards, like other sexually-dimorphic birds, can sometimes go though spontaneous sex reversal, often caused by damaged or nonfunctioning sex organs, such as the ovaries in mallard hens. This phenomenon can cause female mallards to exhibit male plumage, and vice versa (phenotypic feminisation or masculinisation).
Both male and female mallards have distinct iridescent purple-blue speculum feathers edged with white, which are prominent in flight or at rest but temporarily shed during the annual summer moult. Upon hatching, the plumage of the duckling is yellow on the underside and face (with streaks by the eyes) and black on the back (with some yellow spots) all the way to the top and back of the head. Its legs and bill are also black. As it nears a month in age, the duckling's plumage starts becoming drab, looking more like the female, though more streaked, and its legs lose their dark grey colouring. Two months after hatching, the fledgling period has ended, and the duckling is now a juvenile. The duckling is able to fly 50–60 days after hatching. Its bill soon loses its dark grey colouring, and its sex can finally be distinguished visually by three factors: 1)the bill is yellow in males, but black and orange in females; 2)the breast feathers are reddish-brown in males, but brown in females; and 3)in males, the centre tail feather (drake feather) is curled, but in females, the centre tail feather is straight. During the final period of maturity leading up to adulthood (6–10 months of age), the plumage of female juveniles remains the same while the plumage of male juveniles gradually changes to its characteristic colours. This change in plumage also applies to adult mallard males when they transition in and out of their non-breeding eclipse plumage at the beginning and the end of the summer moulting period. The adulthood age for mallards is fourteen months, and the average life expectancy is three years, but they can live to twenty.
Several species of duck have brown-plumaged females that can be confused with the female mallard. The female gadwall (Mareca strepera) has an orange-lined bill, white belly, black and white speculum that is seen as a white square on the wings in flight, and is a smaller bird. More similar to the female mallard in North America are the American black duck (A.rubripes), which is notably darker-hued in both sexes than the mallard, and the mottled duck (A.fulvigula), which is somewhat darker than the female mallard, and with slightly different bare-part colouration and no white edge on the speculum.
In captivity, domestic ducks come in wild-type plumages, white, and other colours. Most of these colour variants are also known in domestic mallards not bred as livestock, but kept as pets, aviary birds, etc., where they are rare but increasing in availability.
A noisy species, the female has the deep quack stereotypically associated with ducks. The female will often call with a sequence of 2–10 quacks in a row, starting loud and with the volume gradually decreasing. Male mallards make a sound phonetically similar to that of the female, a typical quack, but it is deeper and quieter compared to that of the female. Research conducted by Middlesex University on two English mallard populations found that the vocalisations of the mallard varies depending on their environment and have something akin to a regional accent, with urban mallards in London being much louder and more vociferous compared to rural mallards in Cornwall, serving as an adaptation to persistent levels of anthropogenic noise.
When incubating a nest, or when offspring are present, females vocalise differently, making a call that sounds like a truncated version of the usual quack. This maternal vocalisation is highly attractive to their young. The repetition and frequency modulation of these quacks form the auditory basis for species identification in offspring, a process known as acoustic conspecific identification. In addition, females hiss if the nest or offspring are threatened or interfered with. When taking off, the wings of a mallard produce a characteristic faint whistling noise.
The mallard is a rare example of both Allen's Rule and Bergmann's Rule in birds. Bergmann's Rule, which states that polar forms tend to be larger than related ones from warmer climates, has numerous examples in birds, as in case of the Greenland mallard which is larger than the mallards further south. Allen's Rule says that appendages like ears tend to be smaller in polar forms to minimise heat loss, and larger in tropical and desert equivalents to facilitate heat diffusion, and that the polar taxa are stockier overall. Examples of this rule in birds are rare as they lack external ears, but the bill of ducks is supplied with a few blood vessels to prevent heat loss, and, as in the Greenland mallard, the bill is smaller than that of birds farther south, illustrating the rule.
Due to the variability of the mallard's genetic code, which gives it its vast interbreeding capability, mutations in the genes that decide plumage colour are very common and have resulted in a wide variety of hybrids, such as Brewer's duck (mallard × gadwall, Mareca strepera).
Distribution and habitat
The mallard is widely distributed across the Northern and Southern Hemispheres; in North America its range extends from southern and central Alaska to Mexico, the Hawaiian Islands, across the Palearctic, from Iceland and southern Greenland and parts of Morocco (North Africa) in the west, Scandinavia and Britain to the north, and to Siberia, Japan, and South Korea. Also in the east, it ranges to south-eastern and south-western Australia and New Zealand in the Southern hemisphere. It is strongly migratory in the northern parts of its breeding range, and winters farther south. For example, in North America, it winters south to the southern United States and northern Mexico, but also regularly strays into Central America and the Caribbean between September and May. A drake later named "Trevor" attracted media attention in 2018 when it turned up on the island of Niue, an atypical location for mallards.
The mallard inhabits a wide range of habitats and climates, from the Arctic tundra to subtropical regions. It is found in both fresh- and salt-water wetlands, including parks, small ponds, rivers, lakes and estuaries, as well as shallow inlets and open sea within sight of the coastline. Water depths of less than are preferred, with birds avoiding areas more than a few metres deep. They are attracted to bodies of water with aquatic vegetation.
Behaviour
Feeding
The mallard is omnivorous and very flexible in its choice of food. Its diet may vary based on several factors, including the stage of the breeding cycle, short-term variations in available food, nutrient availability, and interspecific and intraspecific competition. The majority of the mallard's diet seems to be made up of gastropods, insects (including beetles, flies, lepidopterans, dragonflies, and caddisflies), crustaceans, other arthropods, worms, feces of other birds, many varieties of seeds and plant matter, and roots and tubers. During the breeding season, male birds were recorded to have eaten 37.6% animal matter and 62.4% plant matter, most notably the grass Echinochloa crus-galli, and nonlaying females ate 37.0% animal matter and 63.0% plant matter, while laying females ate 71.9% animal matter and only 28.1% plant matter. Plants generally make up the larger part of a bird's diet, especially during autumn migration and in the winter.
The mallard usually feeds by dabbling for plant food or grazing; there are reports of it eating frogs, other amphibians, and fish, including carcasses. However, in 2017 a flock of mallards in Romania were observed hunting fledglings of small migratory birds when they land in the water, which included a grey wagtail and a black redstart. This was the first documented occasion they had been seen attacking and consuming large vertebrates. It usually nests on a river bank, but not always near water. It is highly gregarious outside of the breeding season and forms large flocks, which are known as "sordes".
Breeding
Mallards usually form pairs (in October and November in the Northern Hemisphere) until the female lays eggs at the start of the nesting season, which is around the beginning of spring. At this time she is left by the male who joins up with other males to await the moulting period, which begins in June (in the Northern Hemisphere). During the brief time before this, however, the males are still sexually potent and some of them either remain on standby to sire replacement clutches (for female mallards that have lost or abandoned their previous clutch) or forcibly mate with females that appear to be isolated or unattached regardless of their species and whether or not they have a brood of ducklings.
Nesting sites are typically on the ground, hidden in vegetation where the female's speckled plumage serves as effective camouflage, but female mallards have also been known to nest in hollows in trees, boathouses, roof gardens and on balconies, sometimes resulting in hatched offspring having difficulty following their parent to water. Egg clutches number 8–13 creamy white to greenish-buff eggs free of speckles. They measure about in length and in width. The eggs are laid on alternate days, and incubation begins when the clutch is almost complete. Incubation takes 27–28days and fledging takes 50–60days. The ducklings are precocial and fully capable of swimming as soon as they hatch. However, filial imprinting compels them to instinctively stay near the mother, not only for warmth and protection but also to learn about and remember their habitat as well as how and where to forage for food. Though adoptions are known to occur, female mallards typically do not tolerate stray ducklings near their broods, and will violently attack and drive away any unfamiliar young, sometimes going as far as to kill them.
When ducklings mature into flight-capable juveniles, they learn about and remember their traditional migratory routes (unless they are born and raised in captivity). In New Zealand, where mallards are naturalised, the nesting season has been found to be longer, eggs and clutches are larger and nest survival is generally greater compared with mallards in their native range.
In cases where a nest or brood fails, some mallards may mate for a second time in an attempt to raise a second clutch, typically around early-to-mid summer. In addition, mallards may occasionally breed during the autumn in cases of unseasonably warm weather; one such instance of a 'late' clutch occurred in November 2011, in which a female successfully hatched and raised a clutch of eleven ducklings at the London Wetland Centre.
During the breeding season, both male and female mallards can become aggressive, driving off competitors to themselves or their mate by charging at them. Males tend to fight more than females and attack each other by repeatedly pecking at their rival's chest, ripping out feathers and even skin on rare occasions. Female mallards are also known to carry out 'inciting displays', which encourage other ducks in the flock to begin fighting. It is possible that this behaviour allows the female to evaluate the strength of potential partners.
The drakes that end up being left out after the others have paired off with mating partners sometimes target an isolated female duck, even one of a different species, and proceed to chase and peck at her until she weakens, at which point the males take turns copulating with the female. Lebret (1961) calls this behaviour "Attempted Rape Flight", and Stanley Cramp and K.E.L. Simmons (1977) speak of "rape-intent flights". Male mallards also occasionally chase other male ducks of a different species, and even each other, in the same way. In one documented case of "homosexual necrophilia", a male mallard copulated with another male he was chasing after the chased male died upon flying into a glass window. This paper was awarded an Ig Nobel Prize in 2003.
Mallards are opportunistically targeted by brood parasites, occasionally having eggs laid in their nests by redheads, ruddy ducks, lesser scaup, gadwalls, northern shovellers, northern pintails, cinnamon teal, common goldeneyes, and other mallards. These eggs are generally accepted when they resemble the eggs of the host mallard, but the hen may attempt to eject them or even abandon the nest if parasitism occurs during egg laying.
Predators and threats
In addition to human hunting, mallards of all ages (but especially young ones) and in all locations must contend with a wide diversity of predators including raptors and owls, mustelids, corvids, snakes, raccoons, opossums, skunks, turtles, large fish, felids, and canids, the last two including domestic cats and dogs. The most prolific natural predators of adult mallards are red foxes (Vulpes vulpes; which most often pick off brooding females) and the faster or larger birds of prey, (e.g. peregrine falcons, Aquila or Haliaeetus eagles). In North America, adult mallards face no fewer than 15 species of birds of prey, from northern harriers (Circus hudsonius) and short-eared owls (Asio flammeus) (both smaller than a mallard) to huge bald (Haliaeetus leucocephalus) and golden eagles (Aquila chrysaetos), and about a dozen species of mammalian predators, not counting several more avian and mammalian predators who threaten eggs and nestlings.
Mallards are also preyed upon by other waterside apex predators, such as grey herons (Ardea cinerea), great blue herons (Ardea herodias) and black-crowned night herons (Nycticorax nycticorax), the European herring gull (Larus argentatus), the wels catfish (Silurus glanis), and the northern pike (Esox lucius). Crows (Corvus) are also known to kill ducklings and adults on occasion. Also, mallards may be attacked by larger anseriformes such as swans (Cygnus) and geese during the breeding season, and are frequently driven off by these birds over territorial disputes. Mute swans (Cygnus olor) have been known to attack or even kill mallards if they feel that the ducks pose a threat to their offspring. Common loons (Gavia inmer) are similarly territorial and aggressive towards other birds in such disputes, and will frequently drive mallards away from their territory. However, in 2019, a pair of common loons in Wisconsin were observed raising a mallard duckling for several weeks, having seemingly adopted the bird after it had been abandoned by its parents.
In summer, a combination of hot temperatures and reduced water levels place mallards at an increased risk of contracting botulism, as these conditions are ideal for Clostridium botulinum to propagate, with the birds also more likely to come into contact with botulinum toxin produced by the bacteria. Outbreaks of botulism among mallard populations can lead to mass die-offs.
The predation-avoidance behaviour of sleeping with one eye open, allowing one brain hemisphere to remain aware while the other half sleeps, was first demonstrated in mallards, although it is believed to be widespread among birds in general.
Status and conservation
Since 1998, the mallard has been rated as a species of least concern on the IUCN Red List of Endangered Species. This is because it has a large range–more than 20,000,000 km2 (7,700,000 mi2) and because its population is increasing, rather than declining by 30% over ten years or three generations and thus is not warranted a vulnerable rating. Also, the population size of the mallard is very large.
Unlike many waterfowl, mallards have benefited from human alterations to the worldso much so that they are now considered an invasive species in some regions. They are a common sight in urban parks, lakes, ponds, and other human-made water features in the regions they inhabit, and are often tolerated or encouraged in human habitat due to their placid nature towards humans and their beautiful and iridescent colours. While most are not domesticated, mallards are so successful at coexisting in human regions that the main conservation risk they pose comes from the loss of genetic diversity among a region's traditional ducks once humans and mallards colonise an area. Mallards are very adaptable, being able to live and even thrive in urban areas which may have supported more localised, sensitive species of waterfowl before development. The release of feral mallards in areas where they are not native sometimes creates problems through interbreeding with indigenous waterfowl. These non-migratory mallards interbreed with indigenous wild ducks from local populations of closely related species through genetic pollution by producing fertile offspring. Complete hybridisation of various species of wild duck gene pools could result in the extinction of many indigenous waterfowl. The mallard itself is the ancestor of most domestic ducks, and its naturally evolved wild gene pool gets genetically polluted in turn by the domestic and feral populations.
Over time, a continuum of hybrids ranging between almost typical examples of either species develop; the speciation process is beginning to reverse itself. This has created conservation concerns for relatives of the mallard, such as the Hawaiian duck, the New Zealand grey duck (A.s. superciliosa) subspecies of the Pacific black duck, the American black duck, the mottled duck, Meller's duck, the yellow-billed duck, and the Mexican duck, in the latter case even leading to a dispute as to whether these birds should be considered a species (and thus entitled to more conservation research and funding) or included in the mallard species. Ecological changes and hunting have also led to a decline of local species; for example, the New Zealand grey duck population declined drastically due to overhunting in the mid-20th century. Hybrid offspring of Hawaiian ducks seem to be less well adapted to native habitat, and using them in re-introduction projects apparently reduces success. In summary, the problems of mallards "hybridising away" relatives is more a consequence of local ducks declining than of mallards spreading; allopatric speciation and isolating behaviour have produced today's diversity of mallard-like ducks despite the fact that, in most, if not all, of these populations, hybridisation must have occurred to some extent.
Invasiveness
Mallards are causing severe "genetic pollution" to South Africa's biodiversity by breeding with endemic ducks even though the Agreement on the Conservation of African-Eurasian Migratory Waterbirds – an agreement to protect the local waterfowl populations – applies to the mallard as well as other ducks. The hybrids of mallards and the yellow-billed duck are fertile, capable of producing hybrid offspring. If this continues, only hybrids occur and in the long term result in the extinction of various indigenous waterfowl. The mallard can crossbreed with 63 other species, posing a severe threat to indigenous waterfowl's genetic integrity. Mallards and their hybrids compete with indigenous birds for resources, including nest sites, roosting sites, and food.
Availability of mallards, mallard ducklings, and fertilised mallard eggs for public sale and private ownership, either as poultry or as pets, is currently legal in the United States, except for the state of Florida, which has currently banned domestic ownership of mallards. This is to prevent hybridisation with the native mottled duck.
The mallard is considered an invasive species in Australia and New Zealand, where it competes with the Pacific black duck (known as the grey duck locally in New Zealand) which was over-hunted in the past. There, and elsewhere, mallards are spreading with increasing urbanisation and hybridising with local relatives.
The eastern or Chinese spot-billed duck is currently introgressing into the mallard populations of the Primorsky Krai, possibly due to habitat changes from global warming. The Mariana mallard was a resident allopatric population – in most respects a good species – apparently initially derived from mallard-Pacific black duck hybrids; it became extinct in the late 20th century.
The Laysan duck is an insular relative of the mallard, with a very small and fluctuating population. Mallards sometimes arrive on its island home during migration, and can be expected to occasionally have remained and hybridised with Laysan ducks as long as these species have existed. However, these hybrids are less well adapted to the peculiar ecological conditions of Laysan Island than the local ducks, and thus have lower fitness. Laysan ducks were found throughout the Hawaiian archipelago before 400AD, after which they suffered a rapid decline during the Polynesian colonisation. Now, their range includes only Laysan Island. It is one of the successfully translocated birds, after having become nearly extinct in the early 20th century.
Relationship with humans
Domestication
Mallards have often been ubiquitous in their regions among the ponds, rivers, and streams of human parks, farms, and other human-made waterwayseven to the point of visiting water features in human courtyards.
Mallards have had a long relationship with humans. Almost all domestic duck breeds derive from the mallard, with the exception of a few Muscovy breeds, and are listed under the trinomial name A. p. domesticus. Mallards are generally monogamous while domestic ducks are mostly polygamous. Domestic ducks have no territorial behaviour and are less aggressive than mallards. Domestic ducks are mostly kept for meat; their eggs are also eaten, and have a strong flavour. They were first domesticated in Southeast Asia at least 4,000 years ago, during the Neolithic Age, and were also farmed by the Romans in Europe, and the Malays in Asia. As the domestic duck and the mallard are the same species as each other, it is common for mallards to mate with domestic ducks and produce hybrid offspring that are fully fertile. Because of this, mallards have been found to be contaminated with the genes of the domestic duck.
While the keeping of domestic breeds is more popular, pure-bred mallards are sometimes kept for eggs and meat, although they may require wing clipping to restrict flying.
Hunting
Mallards are one of the most common species shot in waterfowl hunting due to their large population size. The ideal location for hunting mallards is considered to be where the water level is somewhat shallow where the birds can be found foraging for food. Hunting mallards might cause the population to decline in some places, at some times, and with some populations. In certain countries, the mallard may be legally shot but is protected under national acts and policies. For example, in the United Kingdom, the mallard is protected under the Wildlife and Countryside Act 1981, which restricts certain hunting methods or taking or killing mallards.
Mallard-Vehicle Collisions
Since standardized data collection began in 1990, the United States Federal Aviation Administration has recorded 1320 mallard collisions with aircraft, 261 of which caused damage to the craft (through 2022). In the United States, the mallard ranks as the 7th most hazardous bird to both military and commercial aircraft. Mallards are of particular concern due to their ubiquity; they are widespread and adaptable to urban environments. Mallards also generally fail to avoid approaching vehicles in experimental settings, especially at high vehicle speeds.
Though most bird strikes occur during the takeoff and landing phases of flight, at least one mallard has been struck at cruising altitude (21,000 feet).
As food
Since ancient times, the mallard has been eaten as food. The wild mallard was eaten in Neolithic Greece. Usually, only the breast and thigh meat is eaten. It does not need to be hung before preparation, and is often braised or roasted, sometimes flavoured with bitter orange or with port.
In culture
Make Way for Ducklings is a children's picture book written and illustrated by Robert McCloskey. The book centers on a pair of mallards who raise their ducklings in the Boston Public Garden.
Migration is an animated adventure comedy film produced by Universal Pictures and Illumination. The story follows a family of mallards who try to migrate from New England to Jamaica.
The world's loneliest duck, named "Trevor" by locals after New Zealand politician Trevor Mallard, appeared without explanation on the Pacific island of Niue, dying there in 2019.
| Biology and health sciences | Anseriformes | null |
230487 | https://en.wikipedia.org/wiki/Poincar%C3%A9%20group | Poincaré group | The Poincaré group, named after Henri Poincaré (1905), was first defined by Hermann Minkowski (1908) as the isometry group of Minkowski spacetime. It is a ten-dimensional non-abelian Lie group that is of importance as a model in our understanding of the most basic fundamentals of physics.
Overview
The Poincaré group consists of all coordinate transformations of Minkowski space that do not change the spacetime interval between events. For example, if everything were postponed by two hours, including the two events and the path you took to go from one to the other, then the time interval between the events recorded by a stopwatch that you carried with you would be the same. Or if everything were shifted five kilometres to the west, or turned 60 degrees to the right, you would also see no change in the interval. It turns out that the proper length of an object is also unaffected by such a shift.
In total, there are ten degrees of freedom for such transformations. They may be thought of as translation through time or space (four degrees, one per dimension); reflection through a plane (three degrees, the freedom in orientation of this plane); or a "boost" in any of the three spatial directions (three degrees). Composition of transformations is the operation of the Poincaré group, with rotations being produced as the composition of an even number of reflections.
In classical physics, the Galilean group is a comparable ten-parameter group that acts on absolute time and space. Instead of boosts, it features shear mappings to relate co-moving frames of reference.
In general relativity, i.e. under the effects of gravity, Poincaré symmetry applies only locally. A treatment of symmetries in general relativity is not in the scope of this article.
Poincaré symmetry
Poincaré symmetry is the full symmetry of special relativity. It includes:
translations (displacements) in time and space, forming the abelian Lie group of spacetime translations (P);
rotations in space, forming the non-abelian Lie group of three-dimensional rotations (J);
boosts, transformations connecting two uniformly moving bodies (K).
The last two symmetries, J and K, together make the Lorentz group (see also Lorentz invariance); the semi-direct product of the spacetime translations group and the Lorentz group then produce the Poincaré group. Objects that are invariant under this group are then said to possess Poincaré invariance or relativistic invariance.
10 generators (in four spacetime dimensions) associated with the Poincaré symmetry, by Noether's theorem, imply 10 conservation laws:
1 for the energy – associated with translations through time
3 for the momentum – associated with translations through spatial dimensions
3 for the angular momentum – associated with rotations between spatial dimensions
3 for a quantity involving the velocity of the center of mass – associated with hyperbolic rotations between each spatial dimension and time
Poincaré group
The Poincaré group is the group of Minkowski spacetime isometries. It is a ten-dimensional noncompact Lie group. The four-dimensional abelian group of spacetime translations is a normal subgroup, while the six-dimensional Lorentz group is also a subgroup, the stabilizer of the origin. The Poincaré group itself is the minimal subgroup of the affine group which includes all translations and Lorentz transformations. More precisely, it is a semidirect product of the spacetime translations group and the Lorentz group,
with group multiplication
.
Another way of putting this is that the Poincaré group is a group extension of the Lorentz group by a vector representation of it; it is sometimes dubbed, informally, as the inhomogeneous Lorentz group. In turn, it can also be obtained as a group contraction of the de Sitter group , as the de Sitter radius goes to infinity.
Its positive energy unitary irreducible representations are indexed by mass (nonnegative number) and spin (integer or half integer) and are associated with particles in quantum mechanics (see Wigner's classification).
In accordance with the Erlangen program, the geometry of Minkowski space is defined by the Poincaré group: Minkowski space is considered as a homogeneous space for the group.
In quantum field theory, the universal cover of the Poincaré group
which may be identified with the double cover
is more important, because representations of are not able to describe fields with spin 1/2; i.e. fermions. Here is the group of complex matrices with unit determinant, isomorphic to the Lorentz-signature spin group .
Poincaré algebra
The Poincaré algebra is the Lie algebra of the Poincaré group. It is a Lie algebra extension of the Lie algebra of the Lorentz group. More specifically, the proper (), orthochronous () part of the Lorentz subgroup (its identity component), , is connected to the identity and is thus provided by the exponentiation of this Lie algebra. In component form, the Poincaré algebra is given by the commutation relations:
where is the generator of translations, is the generator of Lorentz transformations, and is the Minkowski metric (see Sign convention).
The bottom commutation relation is the ("homogeneous") Lorentz group, consisting of rotations, , and boosts, . In this notation, the entire Poincaré algebra is expressible in noncovariant (but more practical) language as
where the bottom line commutator of two boosts is often referred to as a "Wigner rotation". The simplification permits reduction of the Lorentz subalgebra to and efficient treatment of its associated representations. In terms of the physical parameters, we have
The Casimir invariants of this algebra are and where is the Pauli–Lubanski pseudovector; they serve as labels for the representations of the group.
The Poincaré group is the full symmetry group of any relativistic field theory. As a result, all elementary particles fall in representations of this group. These are usually specified by the four-momentum squared of each particle (i.e. its mass squared) and the intrinsic quantum numbers , where is the spin quantum number, is the parity and is the charge-conjugation quantum number. In practice, charge conjugation and parity are violated by many quantum field theories; where this occurs, and are forfeited. Since CPT symmetry is invariant in quantum field theory, a time-reversal quantum number may be constructed from those given.
As a topological space, the group has four connected components: the component of the identity; the time reversed component; the spatial inversion component; and the component which is both time-reversed and spatially inverted.
Other dimensions
The definitions above can be generalized to arbitrary dimensions in a straightforward manner. The -dimensional Poincaré group is analogously defined by the semi-direct product
with the analogous multiplication
.
The Lie algebra retains its form, with indices and now taking values between and . The alternative representation in terms of and has no analogue in higher dimensions.
| Physical sciences | Theory of relativity | Physics |
230488 | https://en.wikipedia.org/wiki/Minkowski%20space | Minkowski space | In physics, Minkowski space (or Minkowski spacetime) () is the main mathematical description of spacetime in the absence of gravitation. It combines inertial space and time manifolds into a four-dimensional model.
The model helps show how a spacetime interval between any two events is independent of the inertial frame of reference in which they are recorded. Mathematician Hermann Minkowski developed it from the work of Hendrik Lorentz, Henri Poincaré, and others said it "was grown on experimental physical grounds".
Minkowski space is closely associated with Einstein's theories of special relativity and general relativity and is the most common mathematical structure by which special relativity is formalized. While the individual components in Euclidean space and time might differ due to length contraction and time dilation, in Minkowski spacetime, all frames of reference will agree on the total interval in spacetime between events. Minkowski space differs from four-dimensional Euclidean space insofar as it treats time differently than the three spatial dimensions.
In 3-dimensional Euclidean space, the isometry group (maps preserving the regular Euclidean distance) is the Euclidean group. It is generated by rotations, reflections and translations. When time is appended as a fourth dimension, the further transformations of translations in time and Lorentz boosts are added, and the group of all these transformations is called the Poincaré group. Minkowski's model follows special relativity, where motion causes time dilation changing the scale applied to the frame in motion and shifts the phase of light.
Spacetime is equipped with an indefinite, non-degenerate, symmetric, bilinear form, called the Minkowski metric, the Minkowski norm squared or Minkowski inner product depending on the context. The Minkowski inner product is defined so as to yield the spacetime interval between two events when given their coordinate difference vector as an argument. Equipped with this inner product (albeit, not technically an inner product), the mathematical model of spacetime is called Minkowski space. The group of transformations for Minkowski space that preserves the spacetime interval (as opposed to the spatial Euclidean distance) is the Poincaré group (as opposed to the Galilean group).
History
Complex Minkowski spacetime
In his second relativity paper in 1905, Henri Poincaré showed how, by taking time to be an imaginary fourth spacetime coordinate , where is the speed of light and is the imaginary unit, Lorentz transformations can be visualized as ordinary rotations of the four-dimensional Euclidean sphere. The four-dimensional spacetime can be visualized as a four-dimensional space, with each point representing an event in spacetime. The Lorentz transformations can then be thought of as rotations in this four-dimensional space, where the rotation axis corresponds to the direction of relative motion between the two observers and the rotation angle is related to their relative velocity.
To understand this concept, one should consider the coordinates of an event in spacetime represented as a four-vector . A Lorentz transformation is represented by a matrix that acts on the four-vector, changing its components. This matrix can be thought of as a rotation matrix in four-dimensional space, which rotates the four-vector around a particular axis.
Rotations in planes spanned by two space unit vectors appear in coordinate space as well as in physical spacetime as Euclidean rotations and are interpreted in the ordinary sense. The "rotation" in a plane spanned by a space unit vector and a time unit vector, while formally still a rotation in coordinate space, is a Lorentz boost in physical spacetime with real inertial coordinates. The analogy with Euclidean rotations is only partial since the radius of the sphere is actually imaginary, which turns rotations into rotations in hyperbolic space (see hyperbolic rotation).
This idea, which was mentioned only briefly by Poincaré, was elaborated by Minkowski in a paper in German published in 1908 called "The Fundamental Equations for Electromagnetic Processes in Moving Bodies". He reformulated Maxwell equations as a symmetrical set of equations in the four variables combined with redefined vector variables for electromagnetic quantities, and he was able to show directly and very simply their invariance under Lorentz transformation. He also made other important contributions and used matrix notation for the first time in this context.
From his reformulation, he concluded that time and space should be treated equally, and so arose his concept of events taking place in a unified four-dimensional spacetime continuum.
Real Minkowski spacetime
In a further development in his 1908 "Space and Time" lecture, Minkowski gave an alternative formulation of this idea that used a real time coordinate instead of an imaginary one, representing the four variables of space and time in the coordinate form in a four-dimensional real vector space. Points in this space correspond to events in spacetime. In this space, there is a defined light-cone associated with each point, and events not on the light cone are classified by their relation to the apex as spacelike or timelike. It is principally this view of spacetime that is current nowadays, although the older view involving imaginary time has also influenced special relativity.
In the English translation of Minkowski's paper, the Minkowski metric, as defined below, is referred to as the line element. The Minkowski inner product below appears unnamed when referring to orthogonality (which he calls normality) of certain vectors, and the Minkowski norm squared is referred to (somewhat cryptically, perhaps this is a translation dependent) as "sum".
Minkowski's principal tool is the Minkowski diagram, and he uses it to define concepts and demonstrate properties of Lorentz transformations (e.g., proper time and length contraction) and to provide geometrical interpretation to the generalization of Newtonian mechanics to relativistic mechanics. For these special topics, see the referenced articles, as the presentation below will be principally confined to the mathematical structure (Minkowski metric and from it derived quantities and the Poincaré group as symmetry group of spacetime) following from the invariance of the spacetime interval on the spacetime manifold as consequences of the postulates of special relativity, not to specific application or derivation of the invariance of the spacetime interval. This structure provides the background setting of all present relativistic theories, barring general relativity for which flat Minkowski spacetime still provides a springboard as curved spacetime is locally Lorentzian.
Minkowski, aware of the fundamental restatement of the theory which he had made, said
Though Minkowski took an important step for physics, Albert Einstein saw its limitation:
For further historical information see references , and .
Causal structure
Where is velocity, , , and are Cartesian coordinates in 3-dimensional space, is the constant representing the universal speed limit, and is time, the four-dimensional vector is classified according to the sign of . A vector is timelike if , spacelike if , and null or lightlike if . This can be expressed in terms of the sign of , also called scalar product, as well, which depends on the signature. The classification of any vector will be the same in all frames of reference that are related by a Lorentz transformation (but not by a general Poincaré transformation because the origin may then be displaced) because of the invariance of the spacetime interval under Lorentz transformation.
The set of all null vectors at an event of Minkowski space constitutes the light cone of that event. Given a timelike vector , there is a worldline of constant velocity associated with it, represented by a straight line in a Minkowski diagram.
Once a direction of time is chosen, timelike and null vectors can be further decomposed into various classes. For timelike vectors, one has
future-directed timelike vectors whose first component is positive (tip of vector located in causal future (also called the absolute future) in the figure) and
past-directed timelike vectors whose first component is negative (causal past (also called the absolute past)).
Null vectors fall into three classes:
the zero vector, whose components in any basis are (origin),
future-directed null vectors whose first component is positive (upper light cone), and
past-directed null vectors whose first component is negative (lower light cone).
Together with spacelike vectors, there are 6 classes in all.
An orthonormal basis for Minkowski space necessarily consists of one timelike and three spacelike unit vectors. If one wishes to work with non-orthonormal bases, it is possible to have other combinations of vectors. For example, one can easily construct a (non-orthonormal) basis consisting entirely of null vectors, called a null basis.
Vector fields are called timelike, spacelike, or null if the associated vectors are timelike, spacelike, or null at each point where the field is defined.
Properties of time-like vectors
Time-like vectors have special importance in the theory of relativity as they correspond to events that are accessible to the observer at (0, 0, 0, 0) with a speed less than that of light. Of most interest are time-like vectors that are similarly directed, i.e. all either in the forward or in the backward cones. Such vectors have several properties not shared by space-like vectors. These arise because both forward and backward cones are convex, whereas the space-like region is not convex.
Scalar product
The scalar product of two time-like vectors and is
Positivity of scalar product: An important property is that the scalar product of two similarly directed time-like vectors is always positive. This can be seen from the reversed Cauchy–Schwarz inequality below. It follows that if the scalar product of two vectors is zero, then one of these, at least, must be space-like. The scalar product of two space-like vectors can be positive or negative as can be seen by considering the product of two space-like vectors having orthogonal spatial components and times either of different or the same signs.
Using the positivity property of time-like vectors, it is easy to verify that a linear sum with positive coefficients of similarly directed time-like vectors is also similarly directed time-like (the sum remains within the light cone because of convexity).
Norm and reversed Cauchy inequality
The norm of a time-like vector is defined as
The reversed Cauchy inequality is another consequence of the convexity of either light cone. For two distinct similarly directed time-like vectors and this inequality is
or algebraically,
From this, the positive property of the scalar product can be seen.
Reversed triangle inequality
For two similarly directed time-like vectors and , the inequality is
where the equality holds when the vectors are linearly dependent.
The proof uses the algebraic definition with the reversed Cauchy inequality:
The result now follows by taking the square root on both sides.
Mathematical structure
It is assumed below that spacetime is endowed with a coordinate system corresponding to an inertial frame. This provides an origin, which is necessary for spacetime to be modeled as a vector space. This addition is not required, and more complex treatments analogous to an affine space can remove the extra structure. However, this is not the introductory convention and is not covered here.
For an overview, Minkowski space is a -dimensional real vector space equipped with a non-degenerate, symmetric bilinear form on the tangent space at each point in spacetime, here simply called the Minkowski inner product, with metric signature either or . The tangent space at each event is a vector space of the same dimension as spacetime, .
Tangent vectors
In practice, one need not be concerned with the tangent spaces. The vector space structure of Minkowski space allows for the canonical identification of vectors in tangent spaces at points (events) with vectors (points, events) in Minkowski space itself. See e.g. or These identifications are routinely done in mathematics. They can be expressed formally in Cartesian coordinates as
with basis vectors in the tangent spaces defined by
Here, and are any two events, and the second basis vector identification is referred to as parallel transport. The first identification is the canonical identification of vectors in the tangent space at any point with vectors in the space itself. The appearance of basis vectors in tangent spaces as first-order differential operators is due to this identification. It is motivated by the observation that a geometrical tangent vector can be associated in a one-to-one manner with a directional derivative operator on the set of smooth functions. This is promoted to a definition of tangent vectors in manifolds not necessarily being embedded in . This definition of tangent vectors is not the only possible one, as ordinary n-tuples can be used as well.
A tangent vector at a point may be defined, here specialized to Cartesian coordinates in Lorentz frames, as column vectors associated to each Lorentz frame related by Lorentz transformation such that the vector in a frame related to some frame by transforms according to . This is the same way in which the coordinates transform. Explicitly,
This definition is equivalent to the definition given above under a canonical isomorphism.
For some purposes, it is desirable to identify tangent vectors at a point with displacement vectors at , which is, of course, admissible by essentially the same canonical identification. The identifications of vectors referred to above in the mathematical setting can correspondingly be found in a more physical and explicitly geometrical setting in . They offer various degrees of sophistication (and rigor) depending on which part of the material one chooses to read.
Metric signature
The metric signature refers to which sign the Minkowski inner product yields when given space (spacelike to be specific, defined further down) and time basis vectors (timelike) as arguments. Further discussion about this theoretically inconsequential but practically necessary choice for purposes of internal consistency and convenience is deferred to the hide box below. | Physical sciences | Theory of relativity | Physics |
230514 | https://en.wikipedia.org/wiki/Intraplate%20earthquake | Intraplate earthquake | An intraplate earthquake occurs in the interior of a tectonic plate, in contrast to an interplate earthquake on the boundary of a tectonic plate. They are relatively rare compared to the more familiar interplate earthquakes. Buildings far from plate boundaries are rarely protected with seismic retrofitting, so large intraplate earthquakes can inflict heavy damage. Examples of damaging intraplate earthquakes are the devastating 2001 Gujarat earthquake, the 2011 Christchurch earthquake, the 2012 Indian Ocean earthquakes, the 2017 Puebla earthquake, the 1811–1812 New Madrid earthquakes, and the 1886 Charleston earthquake. An earthquake that occurs within a subducting plate is known as an intraslab earthquake.
Description
The Earth's crust is made up of seven primary and eight secondary tectonic plates, plus dozens of tertiary microplates. The large plates move very slowly on top of convection currents in the underlying mantle. Because they do not all move in the same direction, plates often directly collide or slide laterally along each other, a tectonic environment that makes interplate earthquakes frequent.
By contrast, relatively few earthquakes occur in intraplate environments away from plate junctures. These earthquakes often occur at the location of ancient failed rifts, partial fractures of existing plates, because they may leave a weakness in the crust vulnerable to regional tectonic strain.
Intraslab earthquakes radiate more seismic energy than interplate earthquakes (megathrust earthquakes) of a similar magnitude. This variation makes seismic energy a better measure for the potential macroseismic effects of an earthquake than the more common seismic moment used to calculate the magnitude .
Examples
Examples of intraplate earthquakes include those in Mineral, Virginia, in 2011 (estimated magnitude 5.8), Newcastle, New South Wales in 1989, New Madrid in 1811 and 1812 (estimated magnitude as high as 8.6), the Boston (Cape Ann) earthquake of 1755 (estimated magnitude 6.0 to 6.3), earthquakes felt in New York City in 1737 and 1884 (both quakes estimated at 5.5 magnitude), and the Charleston earthquake in South Carolina in 1886 (estimated magnitude 6.5 to 7.3). The Charleston quake was particularly surprising because, unlike Boston and New York, the area had almost no history of even minor earthquakes.
In 2001, a large intraplate earthquake devastated the region of Gujarat, India. The earthquake occurred far from any plate boundaries, which meant the region above the epicenter was unprepared for earthquakes. In particular, the Kutch district suffered tremendous damage, where the death toll was over 12,000 and the total death toll was higher than 20,000.
In 2017, the 24–29 km deep magnitude 6.5 Botswana earthquake that shook eastern Botswana occurred at over 300 km from the nearest active plate boundary. The event occurred in an underpopulated area of Botswana.
The 1888 earthquake in Río de la Plata was an intraplate quake, from reactivated faults in the Quilmes Trough, far from the boundaries of the South American plate. With a magnitude greater than 5.0 it was felt "in the cities of Buenos Aires, La Plata and other small towns and villages along the Rio de Plata coastal regions." The towns of Punta del Este and Maldonado in Uruguay were hit by a tsunami generated by the quake.
Causes
Many cities live with the seismic risk of a rare, large intraplate earthquake. The cause of these earthquakes is often uncertain. In many cases, the causative fault is deeply buried and sometimes cannot even be found. Some studies have shown that quakes can be caused by fluids moving up the crust along ancient fault zones. In such circumstances, it is difficult to estimate the seismic hazard for a given city, especially if there was only one earthquake in historical times. Some progress is being made in understanding the fault mechanics driving these earthquakes.
Intraplate earthquakes may be unrelated to ancient fault zones and instead caused by deglaciation or erosion.
Prediction
Scientists continue to search for the causes of these earthquakes, and especially for some indication of how often they recur. The best success has come with detailed micro-seismic monitoring, involving dense arrays of seismometers. In this manner, very small earthquakes associated with a causative fault can be located with great accuracy, and in most cases these line up in patterns consistent with faulting. Cryoseisms can sometimes be mistaken for intraplate earthquakes.
Intraslab earthquake
In seismology, an intraslab earthquake occurs within a subducting plate, known as slabs. They are most frequent in older slabs which are colder, whereas younger slabs that are warmer rarely produces earthquake. They can be detected within these slabs at depths exceeding ; they are also the source of intermediate and deep-focus earthquakes. Intraslab earthquakes at depths are considered shallow earthquakes and can be destructive to cities. One of the deadliest earthquakes of the 20th century was the 1970 Ancash earthquake, measuring 7.9 and occurring off the coast of Peru. The 2001 Nisqually and 1949 Olympia earthquakes were also intraslab events.
| Physical sciences | Seismology | Earth science |
230563 | https://en.wikipedia.org/wiki/Earthquake%20preparedness | Earthquake preparedness | Preparations for earthquakes can consist of survival measures, preparation that will improve survival in the event of an earthquake, or mitigating measures, that seek to minimise the effect of an earthquake. Common survival measures include storing food and water for an emergency, and educating individuals what to do during an earthquake. Mitigating measures can include firmly securing large items of furniture (such as bookcases and large cabinets), TV and computer screens that may otherwise fall over in an earthquake. Likewise, avoiding storing items above beds or sofas reduces the chance of objects falling on individuals.
Planning for a related tsunami, tsunami preparedness, can also be part of earthquake preparedness.
Building design and retrofitting
Seismic codes in earthquake prone areas may have specific requirements designed to increase new buildings' resistance to earthquakes. Older buildings and homes that are not up to code may be modified to increase their resistance. Modification and earthquake resistant design are also employed in elevated freeways and bridges.
Codes are not designed to make buildings earthquake proof in the sense of them suffering zero damage. The goal of most building designs is to reduce earthquake damage to a building such that it protects the lives of occupants and thus tolerance of some limited damage is accepted and considered a necessary trade-off. A supplement or precursor to retrofitting can be the implementation of earthquake-proof furniture.
Earthquake modification techniques and modern building codes are designed to prevent total destruction of buildings for earthquakes of no greater than 8.5 on the Richter Scale. Although the Richter Scale is referenced, the localized shaking intensity is one of the largest factors to be considered in building resiliency.
Types of preparedness
The basic theme behind preparedness is to be ready for an earthquake. Preparedness starts with an individual's everyday life and involves items and training that would be useful in an earthquake. Preparedness continues on a continuum from individual preparedness through family preparedness, community preparedness and then business, non-profit and governmental preparedness. Some organisations blend these various levels. Business continuity planning encourages businesses to have a Disaster Recovery Plan. The US FEMA breaks down preparedness generally into a pyramid, with citizens on the foundational bottom, on top of which rests local government, state government and federal government in that order.
Children may present particular issues and some planning and resources are directly focused on supporting them. The US FEMA has advice noting that "Disasters can leave children feeling frightened, confused, and insecure" whether a child has experienced it first hand, had it happen to a friend or simply seen it on television. People with disabilities or other special needs may have special emergency preparation needs. FEMA's suggestions for people with disabilities include having copies of prescriptions, charging devices for medical devices such as motorized wheel chairs and a week's supply of medication readily available. Preparedness can also cover pets.
Preparedness can also encompass psychological preparedness: resources are designed to support both community members affected by a disaster and the disaster workers serving them.
A multi-hazard approach, where communities are prepared for several hazards, are more resilient than single hazard approaches and have been gaining popularity.
Long term power outages can cause damage beyond the original disaster that can be mitigated with emergency generators or other power sources to provide an emergency power system. The United States Department of Energy states: "homeowners, business owners, and local leaders may have to take an active role in dealing with energy disruptions on their own." Major institutions like hospitals, military bases and educational institutions often have extensive backup power systems.
Preparedness does not stop at home or at school. The United States Department of Health and Human Services addresses specific emergency preparedness issues hospitals may have to respond to, including maintaining a safe temperature, providing adequate electricity for life support systems and even carrying out evacuations under extreme circumstances. FEMA encourages all businesses to have an emergency response plan and the Small Business Administration specifically advises small business owners to also focus emergency preparedness and provides a variety of different worksheets and resources.
Given the explosive danger posed by natural gas leaks, Ready.gov states that "It is vital that all household members know how to shut off natural gas" and that property owners must ensure they have any special tools needed for their particular gas connections. Ready.gov also notes that "It is wise to teach all responsible household members where and how to shut off the electricity," cautioning that individual circuits should be shut off before the main circuit. Ready.gov further states that "It is vital that all household members learn how to shut off the water at the main house valve" and cautions that the possibility that rusty valves might require replacement.
Achieving preparedness
Levels of preparedness generally remain low, despite attempts to increase public awareness.
Various methods exist to promote disaster preparedness, but they are rarely well documented and their efficacy is rarely tested. Hands on training, drills and face-to-face interaction have proven more successful at changing behaviour. Digital methods have also been used, including for examples educational videogames.
| Physical sciences | Seismology | Earth science |
230613 | https://en.wikipedia.org/wiki/Scythe | Scythe | A scythe (, rhyming with writhe) is an agricultural hand tool for mowing grass or harvesting crops. It is historically used to cut down or reap edible grains, before the process of threshing. The scythe has been largely replaced by horse-drawn and then tractor machinery, but is still used in some areas of Europe and Asia. Reapers are bladed machines that automate the cutting of the scythe, and sometimes subsequent steps in preparing the grain or the straw or hay.
The word "scythe" derives from Old English siðe. In Middle English and later, it was usually spelled sithe or sythe. However, in the 15th century some writers began to use the sc- spelling as they thought (wrongly) the word was related to the Latin scindere (meaning "to cut"). Nevertheless, the sithe spelling lingered and notably appears in Noah Webster's dictionaries.
A scythe consists of a shaft about long called a snaith, snath, snathe or sned, traditionally made of wood but now sometimes metal. Simple snaiths are straight with offset handles, others have an "S" curve or are steam bent in three dimensions to place the handles in an ergonomic configuration but close to the shaft. The snaith has either one or two short handles at right angles to it, usually one near the upper end and always another roughly in the middle. The handles are usually adjustable to suit the user. A curved, steel blade between long is mounted at the lower end at 90°, or less, to the snaith. Scythes almost always have the blade projecting from the left side of the snaith when in use, with the edge towards the mower; left-handed scythes are made but cannot be used together with right-handed scythes as the left-handed mower would be mowing in the opposite direction and could not mow in a team. Although left-handed scythes exist, their primary purpose is not to suit left-handed mowers but to mow back out from an obstruction on the left, such as when mowing back from the end of a ditch; ditch mowers may have both left and right-handed ditch scythes with them to do this.
Use
The use of a scythe was historically called mowing, now often scything to distinguish it from machine mowing. The mower holds the top handle in the left hand and the central one in the right, with the arms straight, the blade parallel and very close to the ground. The mower cuts along the mowing edge of the meadow, keeping the uncut grass to the right. The blade hooks the grass on the right and is swung to the left in a long arc ending to the left of the mower to form a windrow of cut grass on the previously mown ground. The mower takes a small step forward and repeats the motion, proceeding with a steady rhythm, stopping at frequent intervals to hone the blade. The correct technique has a slicing action on the grass, leaving a swathe of uniformly cut stubble, and forming a regular windrow on the left.
When mown in a team, the team starts at the edge of the meadow in a staggered line, then proceeds clockwise, finishing in the middle. Mowing with a scythe is a skilled task that takes time to learn fully. Long-bladed scythes, typically around (such as in the example below) and suitable for mowing grass or wheat, are harder to use at first; consequently, beginners usually start on shorter blades, generally or less. Common beginner errors include setting up the snaith with the handles in the wrong locations to suit the body, setting the blade at the wrong turn-in and turn-up angles to suit the conditions, choosing a blade that is too long for the skill level, failing to start with a sharp edge and persevering with a dull one during use, chopping or hacking at the grass, trying to cut too wide a strip of grass at once and striking the ground with the blade. Historically, beginners relied on mentors to help them set up and maintain their scythe and to teach them to mow comfortably and efficiently.
Mowing grass is easier when it is damp, and so hay-making historically began at dawn and often stopped early, the heat of the day being spent raking and carting the hay cut on previous days or peening the blades.
Scythes are designed for different tasks. A long, thin blade is most efficient for mowing grass or wheat, while a shorter, more robust scythe is more appropriate for clearing weeds, cutting reed or sedge and can be used with the blade under water for clearing ditches and waterways. Skilled mowers using long-bladed scythes honed very sharp were used to maintain short lawn grass until the invention of the lawnmower. Many cultures have used a variety of 'cradles' to catch cut different kinds of grain stems, keeping the seed heads aligned and laying them down in an orderly fashion to make them easier to sheaf and winnow.
The mowing action
Sharpening
The cutting edge of a scythe blade is maintained by occasional peening followed by frequent honing. Peening reforms the malleable edge by hammering; creating the desired edge profile, locally work-hardening the metal, and removing minor nicks and dents. For mowing fine grass, the bevel angle may be peened extremely fine, while for coarser work a larger angle is created to give a more robust edge. Peening requires some skill and is done using a peening hammer and special anvils or by using a peening jig. Historically, when mowing in teams, a peening station was set up on the edge of the field during harvest, but now more likely back in the workshop.
In the example below, a short scythe blade, being used to clear brambles, is being sharpened. Before being taken to the forest, the blade is peened in the workshop: this reforms the malleable steel to create an edge profile that can then be honed. Peening is done only occasionally; how often depends on the hardness of the steel and the nature of the work. The Austrian blade shown is being used to cut tough-stemmed brambles and it is being peened about every thirty hours of work. Nicks and cuts to the blade edge can usually be worked out of the blade by peening and a new edge profile formed for honing.
A peening jig is being used here, but blades can be free-peened using various designs of peening anvils. The peening jig shown has two interchangeable caps that set different angles: a coarse angle is set first about 3 mm back from the edge, and the fine angle is then set on the edge, leaving an edge that lends itself to being easily honed. The blade is then honed using progressively finer honing stones and then taken to the field. In the field, the blade is honed using a fine, ovoid whetstone (or rubber), fine-grained for grass, coarser for cereal crops. Honing is performed the moment the mower senses that the edge has gone off; this may be every half hour or more depending on the conditions. The laminated honing stone shown here has two grades of stone and is carried into the field soaking in a water-filled holster on the belt. A burr is set up on the outside of the blade by stroking the blade on the inside; the burr is then taken off by gently stroking it on the outside. There are many opinions, regional traditions and variations on exactly how to do this; some eastern European countries even set up the burr on the inside.
Unlike European blades, which are made from malleable steel, typical American blades are made of harder, more brittle, steel and risk cracking if peened. While the harder blade holds an edge longer and requires less frequent honing in the field, the edge can only be reshaped by grinding after heavy use or damage. This usually only needs to be done only 1–3 times a season because of the greater wear resistance of the harder steel. Some examples have a laminated construction with a hard, wear-resistant core providing the edge, and softer sides providing strength. In American blades, the edge steel is typically clad on either side with the tough iron, while some Nordic laminated blades have a layer of iron on the top only, with the edge steel comprising the bottom layer.
History
Scythes may date back as far as ; they seem to have been used since Cucuteni–Trypillia settlements, becoming widespread with agricultural developments. Initially used mostly for mowing hay, it had replaced the sickle for reaping crops by the 16th century, as the scythe was better ergonomically and consequently more efficient. In about 1800, the grain cradle was sometimes added to the standard scythe when mowing grain; the cradle was an addition of light wooden fingers above the scythe blade which kept the grain stems aligned and the heads together to make the collection and threshing easier. In the developed world, the scythe has been largely replaced by the motorised lawn mower and combine harvester. However, the scythe remained in common use for many years after the introduction of machines because a side-mounted finger-bar mowerwhether horse- or tractor-drawncould not mow in front of itself, and scythes were still needed to open up a meadow by clearing the first swathe to give the mechanical mower room to start.
The Dictionary of Greek and Roman Antiquities of Sir William Smith argues that the scythe, known in Latin as the falx foenaria as opposed to the sickle, the falx messoria, was used by the ancient Romans. According to ancient Greek mythology, Gaiathe Greek goddess and mother of the Titansgave a sickle made of the strongest metal to her youngest son Kronos, who is also the youngest of the Titans and god of the harvest, to seek vengeance against her husband Ouranos for torturing their eldest sons. The Grim Reaper is often depicted carrying or wielding a scythe. According to Jack Herer and Flesh of The Gods (Emboden, W. A. Jr., Praeger Press, New York, 1974), the ancient Scythians grew hemp and harvested it with a hand reaper that would be considered a scythe.
The Abbeydale Industrial Hamlet in Sheffield, England, is a museum of a scythe-making works that was in operation from the end of the 18th century until the 1930s. This was part of the former scythe-making district of north Derbyshire, which extended into Eckington. Other English scythe-making districts include that around Belbroughton.
The German Renaissance scythe sword, the Greek and Roman harpe and the Egyptian khopesh were scythes or sickles modified as weapons or symbols of authority. An improvised conversion of the agricultural scythe to a war scythe by re-attaching the blade parallel to the snaith, similar to a bill, has also been used throughout history as a weapon. See below for an example.
In national cultures
Although the scythe is still an indispensable tool for farmers in developing countries and in mountainous terrain, the tool has become associated with the Grim Reaper.
In Romania, for example, in the highland landscape of the Transylvanian Apuseni Mountains, scything is a very important annual activity, taking about 2–3 weeks to complete for a regular house. As scything is a tiring physical activity and is relatively difficult to learn, farmers help each other by forming teams. After each day's harvest, the farmers often celebrate by having a small feast where they dance, drink and eat, while being careful to keep in shape for the next day's hard work. In other parts of the Balkans, such as in Serbian towns, scything competitions are held where the winner takes away a small silver scythe. In small Serbian towns, scything is treasured as part of the local folklore, and the winners of friendly competitions are rewarded richly with food and drink, which they share with their competitors.
Among Basques scythe-mowing competitions are still a popular traditional sport, called segalaritza (from Spanish verb segar: to mow). Each contender competes to cut a defined section of grown grass before his rival does the same.
There is an international scything competition held at Goričko where people from Austria, Hungary, Serbia, and Romania, or as far away as Asia enter to showcase their culturally unique method of reaping crops. In 2009, a Japanese man showcased a wooden reaping tool with a metal edge, which he used to show how rice was cut. He was impressed with the speed of the local reapers, but said such a large scythe would never work in Japan.
The Norwegian municipality of Hornindal has three scythe blades in its coat-of-arms.
Their motto was They Feed and Defend (Polish: Żywią i Bronią, archaic spelling: Żywią y Bronią).
The emblem of Bnei Akiva, a Jewish religious-Zionist youth movement, contains wheat, scythe and pitchfork, representing agriculture and the combination of land labor with the Torah.
In 2012, The Wall Street Journal reported some American homeowners were eschewing motorized lawn mowers in favor of scythes, citing the lack of noise as well as the health benefits of scythe reaping as exercise.
In art
| Technology | Agricultural tools | null |
230643 | https://en.wikipedia.org/wiki/Braid | Braid | A braid (also referred to as a plait; ) is a complex structure or pattern formed by interlacing three or more strands of flexible material such as textile yarns, wire, or hair.
The simplest and most common version is a flat, solid, three-stranded structure. More complex patterns can be constructed from an arbitrary number of strands to create a wider range of structures (such as a fishtail braid, a five-stranded braid, rope braid, a French braid and a waterfall braid). The structure is usually long and narrow with each component strand functionally equivalent in zigzagging forward through the overlapping mass of the others. It can be compared with the process of weaving, which usually involves two separate perpendicular groups of strands (warp and weft).
Historically, the materials used have depended on the indigenous plants and animals available in the local area. During the Industrial Revolution, mechanized braiding equipment was invented to increase production. The braiding technique was used to make ropes with both natural and synthetic fibers as well as coaxial cables for radios using copper wire. In more recent times it has been used to create a covering for fuel pipes in jet aircraft and ships (first using glass fibre, then stainless steel and Kevlar). Hoses for domestic plumbing are often covered with stainless steel braid.
Hair braiding
The oldest known reproduction of hair braiding may go back about 30,000 years: the Venus of Willendorf, a female figurine estimated to have been made between about 28,000 and 25,000 BC in modern-day Austria. The Venus of Brassempouy from the southwest of France is estimated to be about 25,000 years old and shows a braided hairstyle.
Although many cultures want to take sole credit for the braid, they cannot be traced to a single origin. Like how different versions of Cinderella are traceable to nearly every culture, braids, too, are polygenetic. One early example of hair braiding takes place in 1279-1213 BCE as recorded in the story of Isis: "when some of the queen's maidens came to the well, she greeted them kindly and began to braid their hair."
During the Bronze Age and Iron Age, many peoples in the Near East, Asia Minor, Caucasus, East Mediterranean and North Africa are depicted in art with braided or plaited hair and beards. Similarly, the practice is recorded in Europe, Africa, India, China, Japan, Australasia and Central Asia.
Braiding is traditionally a social art. Because of the time it takes to braid hair, people have often taken time to socialize while braiding and having their hair braided. It begins with the elders making simple knots and braids for younger children. Older children watch and learn from them, start practicing on younger children, and eventually learn the traditional designs. This carries on a tradition of bonding between elders and the new generation.
Industrial history and use
Early braids had many uses, such as costume decoration, animal regalia (like camel girths), sword decoration, bowls and hats (from palm leaves), locks (such as those made in Japan to secure precious tea supplies through the use of elaborate knots), and weapons (e.g. slings).
Materials that are used in braids can vary depending on local materials. For instance, South Americans used the very fine fibers from the wool of alpaca and llama, while North American people made use of bison fibers. Throughout the world, vegetable fibers such as grass, nettle, and hemp have been used to create braids. In China, Korea, and Japan silk still remains the main material used. In the Americas, the braiding of leather is also common. Plaiting with kangaroo leather has been a widely practiced tradition in rural Australia since pioneering times. It is used in the production of fine leather belts, hatbands, bridles, dog leads, bullwhips, stockwhips, etc. Other leathers are used for the plaiting of heavier products suitable for everyday use.
For nomadic peoples, braiding was a practical means of producing useful and decorative textiles. In other areas, such as the Pacific islands (where leaves and grasses are braided), and for many hill tribes, braids are made using minimal equipment. It was only when braiding became a popular occupation in the home or school, as it is in China and Japan, and when the Industrial Revolution came about, that specific tools were developed to increase production and make it easier to produce more complicated patterns of braids.
Braids are also very good for making rope and decorative objects. Complex braids have been used to create hanging fibre artworks.
Gold braids and silver braids are components or trims of many kinds of formal dress, including military uniform (in epaulettes, aiguillettes, on headgear).
Ropes and cables
Braiding creates a composite rope that is thicker than the non-interlaced strands of yarns. Braided ropes are preferred by arborists, rock climbers, and in sport sailing because they do not twist under load, as does an ordinary twisted-strand rope. These ropes consist of one or more concentric tubular braided jackets surrounding either several small twisted fibre cords, or a single untwisted yarn of straight fibres, and are known as Kernmantle ropes.
In electrical and electronic cables, braid is a tubular sheath made of braided strands of metal placed around a central cable for shielding against electromagnetic interference. The braid is grounded while the central conductor(s) carries the signal. The braid may be used in addition to a foil jacket to increase shielding and durability. Litz wire uses braids of thin insulated wires to carry high frequency signals with much lower losses from skin effect or to minimise proximity effect in transformers. Flat braids made of many copper wires can also be used for flexible electrical connections between large components. The numerous smaller wires comprising the braid are much more resistant to breaking under repeated motion and vibration than is a cable of larger wires. A common example of this may be found connecting a car battery's negative terminal to the metal chassis.
Similar braiding is used on pressurized hoses, such as in plumbing and hydraulic brake systems in automobiles. Braiding is also used for fibres for composite reinforcements.
A property of the basic braid is that removing one strand unlinks the other two, as they are not twisted around each other. Mathematically, a braid with that property is called a Brunnian braid.
Onion and garlic
Onion and garlic stalks are often braided for storage after they are partially dried.
Metaphors
Braids are often used figuratively to represent interweaving or combination, such as in, "He braided many different ideas into a new whole."
In some river and stream systems, small streams join and redivide in many places. Such stream systems are said to be braided. These are often found in alluvial fans at the outlet of canyons. This is a result of heavy sediment deposition at high flows followed by re-erosion at low flows.
Gallery
| Technology | Weaving | null |
230909 | https://en.wikipedia.org/wiki/Whitefly | Whitefly | Whiteflies are Hemipterans that typically feed on the undersides of plant leaves. They comprise the family Aleyrodidae, the only family in the superfamily Aleyrodoidea. More than 1550 species have been described.
Description and taxonomy
The Aleyrodidae are a family in the suborder Sternorrhyncha and at present comprise the entire superfamily Aleyrodoidea, related to the superfamily Psylloidea. The family often occurs in older literature as "Aleurodidae", but that is a junior synonym and accordingly incorrect in terms of the international standards for zoological nomenclature.
Aleyrodidae are small insects, most species with a wingspan of less than 3 mm and a body length of 1 mm to 2 mm. Many are so small that their size complicates their control in greenhouses because they can only be excluded by screening with very fine mesh; in fact they can enter mesh so fine that many of their natural enemies cannot come in after them, so that unchecked whitefly populations in greenhouses rapidly become overwhelming. Some "giant whitefly" species exist, some of which may exceed 5 mm in size. This sometimes is associated with sexual dimorphism in which one sex is markedly larger than the other. Such dimorphism is common in the Sternorrhyncha, in which the males of most scale insects for example are tiny compared to the female. Remarkably however, in some giant tropical species the males are much larger than the females.
Like most of the mobile Sternorrhyncha, adult Aleyrodidae have well-developed antennae, which in most species in this family are seven-segmented.
As in many Hemiptera, there are two ocelli, which generally in the Aleyrodidae are placed at the anterior margins of the compound eyes. The compound eyes themselves are rather remarkable: many have a distinct constriction between the upper and lower halves, and in some species there is a complete separation. Many insects' compound eyes are divided into functionally and anatomically distinct upper and lower regions, but the adaptation's purpose or origin in Aleyrodidae is unclear. The degree of separation is useful in recognising the species; for instance, one way to tell adult Bemisia from Trialeurodes is that the upper and lower parts of the compound eyes are connected by a single ommatidium in Bemisia, while in Trialeurodes they are completely separate.
Both sexes have functional mouthparts and two pairs of membranous, functional wings; the rear wings are neither much reduced, nor modified into any such hooked or haltere-like structures as occur in some other Hemiptera such as many of the Coccoidea. The wing venation is reduced, like that of the Psyllidae, only generally much more so. In many genera there is only one conspicuous and unbranched vein in each wing; however, wings of larger species such as Udamoselis have less reduced venation, though their veins still are simple and few.
The insects and their wings are variously marked or mottled according to species, and many species are covered with fine wax powder, giving most species a floury, dusted appearance, hence names such as Aleyrodidae, Aleurodidae and Aleuroduplidens; the root refers to the (aleurodes) meaning "floury". However, not all species are white; for example, Aleurocanthus woglumi is slaty black.
The legs of Aleyrodidae are well developed and fairly long, but gracile, and in contrast to Psyllidae, not adapted to leaping. The tarsi have two segments of roughly equal length. The pretarsus has paired claws, with an empodium between—in some species the empodium is a bristle, but in others it is a pad.
The digestive system of the Aleyrodidae is typical of the Sternorrhyncha, including a filter chamber, and all active stages of the Aleyrodidae accordingly produce large quantities of honeydew; the anus is adapted to presentation of honeydew to symbiotic species, mainly ants; the honeydew emerges from the anus, which is inside an opening called the vasiform orifice on the dorsal surface of the caudal segment of the abdomen. This orifice is large and is covered by an operculum. The entire structure is characteristic of the Aleyrodidae and within the family it is taxonomically diagnostic because it varies in shape according to the species. Within the orifice beneath the operculum there is a tongue-like lingula. It appears to be involved in the expulsion of honeydew, and in fact at one time was wrongly assumed to be the organ that produced the honeydew. In some species it protrudes from beneath the operculum, but in others it normally is hidden.
Evolutionary history
The oldest members of the family belong to the Mesozoic subfamily Bernaeinae, known from the Middle/Upper Jurassic-Upper Cretaceous, the oldest representatives of the extant subfamilies Aleyrodinae and Aleurodicinae appear during the Lower Cretaceous.
Reproduction and metamorphosis
The eggs of Aleyrodidae generally are laid near each other on the food plant, usually on a leaf, in spiral patterns or arcs, sometimes in parallel arcs. The egg is elongated, with one narrow end produced into a pedicel, which in some species is longer than the rest of the egg. After fertilisation the pedicel shrivels into a stalk.
The details vary, but at least some species can reproduce parthenogenically by automixis. However, apparently all males are parthenogenically produced by arrhenotoky. The female however, can mate with her own male offspring, and thereafter produce eggs of both sexes.
There generally are four larval instars. All the instars are more or less in the shape of a flattened ellipse fringed with bristles and waxy filaments. The first instar has functional legs, though short. Once it has inserted its stylets into the phloem to feed, it settles down and no longer uses its legs, and they degenerate after the first ecdysis. From then until it emerges as an adult, it remains attached to the plant by its mouthparts. The final instar feeds for a while, then undergoes changes within its skin, ceasing feeding and growing a new skin, forming what amounts to pupa. In doing so the insect does not shed the larval skin, which it retains as a protective puparium and which dries out. Meanwhile, the pupa within this skin develops into a pharate adult that usually is visible through the wall of the puparium. The puparium splits open as the imago forces its way out.
This pupal stage is analogous to the pupal forms of the Holometabola and it raises questions of terminology and concept. Some authorities argue that there is little functional, and no logically cogent basis for the distinction between the terms "larva" and "nymph". Some have long been in favour of dropping the term nymph entirely, and certainly apply the term "larvae" to the Aleyrodidae.
Agricultural threat
In warm or tropical climates and especially in greenhouses, whiteflies present major problems in crop protection. Worldwide economic losses are estimated at hundreds of millions of dollars annually.
Prominent pest species include:
Aleurocanthus woglumi, citrus blackfly, which, in spite of its color, is a whitefly that attacks citrus
Aleyrodes proletella, cabbage whitefly, is a pest of various Brassica crops.
Bemisia tabaci, silverleaf whitefly, is a pest of many agricultural and ornamental crops.
Trialeurodes vaporariorum, greenhouse whitefly, a major pest of greenhouse fruit, vegetables, and ornamentals
Although several species of whitefly may cause some crop losses simply by sucking sap when they are very numerous, the major harm they do is indirect. Firstly, like many other sap-sucking Hemiptera, they secrete large amounts of honeydew that support unsightly or harmful infestations of sooty mold. Recent studies suggest that insecticides can also be excreted through the honeydew leading to unintended effects. Secondly, they inject saliva that may harm the plant more than either the mechanical damage of feeding or the growth of the fungi. However, by far their major importance as crop pests is their transmission of diseases of plants.
Although there are a great many species of whiteflies, and the family is notorious for devastating transmission of crop viruses, the actual proportion of whiteflies which are responsible is very low. The most prominent disease vectors among the Aleyrodidae are a species complex in the genus Bemisia. Bemisia tabaci and B. argentifolii transmit African cassava mosaic, bean golden mosaic, bean dwarf mosaic, bean calico mosaic, tomato yellow leaf curl, tomato mottle, and other Begomoviruses, in the family Geminiviridae. The worldwide spread of emerging biotypes, such as B. tabaci biotype B, also known as, 'B. argentifolii', and a new biotype Q, continue to cause severe crop losses which are expected to increase, demanding matching increases in pesticide use on many crops (tomatoes, beans, cassava, cotton, cucurbits, potatoes, sweet potatoes). Efforts to develop environmentally friendly integrated pest management systems, with the goal of reducing insecticide use, aim to re-establish the ecological equilibrium of predators, parasitoids, and microbial controls that were once in place. New crop varieties are also being developed with increased tolerance to whiteflies, and to the plant diseases carried by them. A major problem is that whiteflies and the viruses they carry can infect many host plants, including agricultural crops, palms, and weeds. These problems are complicated by difficulties in classifying and detecting new whitefly biotypes and begomoviruses. Proper diagnosis of plant diseases depends on using sophisticated molecular techniques to detect and characterize the viruses and whiteflies which are present in a crop. A team of researchers, extension agents and growers working together are needed to follow disease development, using dynamic modeling, to understand the incidence of disease spread.
In 1997, tomato yellow leaf-curl begomovirus was discovered in Florida, USA. This is the worst viral disease transmitted by the whitefly, Bemisia argentifolii. The whitefly has also been shown to transmit almost 60 other viral plant diseases.
In 2023 flower plantations in Naga, Camarines Sur were infested by whiteflies. Eggplant crops in Baler, Aurora were also affected. In 2024, the locos Norte City Agriculture Office discovered cauliflower Aleyrodidae infestation in Laoag due to changing weather.
Damage by feeding
Whiteflies feed by tapping into the phloem of plants, introducing toxic saliva and decreasing the plants' overall turgor pressure. Since whiteflies congregate in large numbers, susceptible plants can be quickly overwhelmed. Further harm is done by mold growth encouraged by the honeydew whiteflies secrete. This may also seriously impede the ability of farms to process cotton harvests.
Whiteflies share a modified form of hemimetabolous metamorphosis, in that the immature stages begin life as mobile individuals, but soon attach to host plants. The stage before the adult is called a pupa, though it shares little in common with the pupal stage of holometabolous insects.
Control
Whitefly control is difficult and complex, as whiteflies rapidly develop resistance to chemical pesticides. The USDA recommends "an integrated program that focuses on prevention and relies on cultural and biological control methods when possible".
While an initial pesticide application may be necessary to control heavy infestations, repeated applications may lead to strains of whiteflies that are resistant to pesticides, so only use of selective insecticides is advised. Specific insecticide information and guidance for the fig whitefly is available from the University of Florida. Care should be taken to ensure that the insecticide used will not kill the natural predators of whiteflies. For effective use of biological method after application of pesticide, plant washing is advised prior to release of predators or parasitoids.
Pesticides used for whitefly control usually contain neonicotinoid compounds as active ingredients: clothianidin (commercial), dinotefuran (over-the-counter and commercial), imidacloprid (over-the-counter and commercial) and thiamethoxam (commercial). Neonicotinoids can be harmful if ingested. Rotation of insecticides from different families may be effective at preventing the building of tolerance to the product. Clothianidin and dinotefuran are of the same family. Spraying the leaves using insecticidal soap is another, environmentally friendly, option.
Nonchemical means
Biological methods have also been proposed to control whitefly infestation, and may be paired with chemical methods. Washing the plant, especially the undersides of leaves, may help reduce the number of the pests on the plants and make their management by other methods more effective. Whiteflies are also attracted by the color yellow, so yellow sticky paper can serve as traps to monitor infestations. Dead leaves or leaves that have been mostly eaten by whiteflies can be removed and burned or carefully placed in closed bins to avoid reinfestation and spreading of the disease.
Early detection in combination with hosing or vacuuming of diseased portions as well as removal of any section that is heavily infested. Pesticide use is not ideal in the case of controlling whitefly and widespread contamination can be costly; it is best to avoid this problem with aggressive preventive measures.
Several predators and parasitoids may be effective in controlling whitefly infestations, including green lacewings, ladybirds, minute pirate bugs, big-eyed bugs, damsel bugs, Encarsia formosa and phytoseiid mites.
Integrated management of whiteflies can as well be done using biopesticides based on microbials such as Beauveria bassiana (effective on larvae and adults) or Isaria fumosorosea.
Green lacewing larvae have voracious appetites, so will attack whiteflies, as well as other pests, including aphids, mealybugs, spider mites, leafhopper larvae, moth eggs, scales, and thrips. They also can attack other insects, including caterpillars. They are available as eggs from commercial insectaries, and will stay in a larval stage after they hatch for one to three weeks. The adult insects can fly and will feed only on pollen, honey, and nectar to reproduce. Repeated application may be necessary and the eggs could be eaten before they hatch by their natural predators, such as ants or mature green lacewings.
Ladybirds are also used. They eat mostly insect eggs, but will also feed on beetle larvae, aphids, scale insects, and young caterpillars. Adults are often collected when in a dormant state in the wild and shipped for use in pest control; however, they may not stay in the location where they are released. They do live for about a year and will continuously lay eggs and reproduce. Spraying the bugs' wings with a sticky substance before release may hinder their ability to fly.
Some promising claims have been made that mesh or film that excludes ultraviolet of certain wavelengths from a greenhouse interfere severely with the ability of whitefly and various other greenhouse pests, to find their food plants. It is not yet clear, assuming that the effect is substantially of value, how readily pests in such circumstances might develop behavioural tolerance to such control measures.
Companion plants
A number of plants can be intercropped with vegetables, in a garden setting, serving as companion plants to protect against whiteflies.
For example, nasturtiums are thought to provide a defense to gooseberries or tomatoes. They provide root chemicals that deter whiteflies.
A study intercropping French marigolds with short-vine tomatoes in glasshouse growing conditions achieved some control of whiteflies when the plants were grown together from the beginning of the growing season; however, limonene dispensers were more effective.
Zinnias, conversely, attract predators that consume whiteflies, including hummingbirds and predatory wasps and flies. Other plants with a similar function include the hummingbird bush, pineapple sage, and bee balm. Each of these plants also conceals the scent of nearby plants, making their detection by some pest insects more difficult, as do most other mints.
| Biology and health sciences | Hemiptera (true bugs) | null |
230920 | https://en.wikipedia.org/wiki/Beluga%20whale | Beluga whale | The beluga whale (; Delphinapterus leucas) is an Arctic and sub-Arctic cetacean. It is one of two members of the family Monodontidae, along with the narwhal, and the only member of the genus Delphinapterus. It is also known as the white whale, as it is the only cetacean to regularly occur with this colour; the sea canary, due to its high-pitched calls; and the melonhead, though that more commonly refers to the melon-headed whale, which is an oceanic dolphin.
The beluga is adapted to life in the Arctic, with anatomical and physiological characteristics that differentiate it from other cetaceans. Amongst these are its all-white colour and the absence of a dorsal fin, which allows it to swim under ice with ease. It possesses a distinctive protuberance at the front of its head which houses an echolocation organ called the melon, which in this species is large and deformable. The beluga's body size is between that of a dolphin and a true whale, with males growing up to long and weighing up to . This whale has a stocky body. Like many cetaceans, a large percentage of its weight is blubber (subcutaneous fat). Its sense of hearing is highly developed and its echolocation allows it to move about and find breathing holes under sheet ice.
Belugas are gregarious and form groups of 10 animals on average, although during the summer, they can gather in the hundreds or even thousands in estuaries and shallow coastal areas. They are slow swimmers, but can dive to below the surface. They are opportunistic feeders and their diets vary according to their locations and the season. The majority of belugas live in the Arctic Ocean and the seas and coasts around North America, Russia, and Greenland; their worldwide population is thought to number around 200,000. They are migratory and the majority of groups spend the winter around the Arctic ice cap; when the sea ice melts in summer, they move to warmer river estuaries and coastal areas. Some populations are sedentary and do not migrate over great distances during the year.
The native peoples of North America and Russia have hunted belugas for many centuries. They were also hunted by non-natives during the 19th century and part of the 20th century. Hunting of belugas is not controlled by the International Whaling Commission, and each country has developed its own regulations in different years. Currently, some Inuit in Canada and Greenland, Alaska Native groups and Russians are allowed to hunt belugas for consumption as well as for sale, as aboriginal whaling is excluded from the International Whaling Commission 1986 moratorium on hunting. The numbers have dropped substantially in Russia and Greenland, but not in Alaska and Canada. Other threats include natural predators (polar bears and killer whales), contamination of rivers (as with polychlorinated biphenyl (PCBs) which bioaccumulate up the food chain), climate change and infectious diseases. The beluga was placed on the International Union for Conservation of Nature's Red List in 2008 as being "near threatened"; the subpopulation from the Cook Inlet in Alaska is considered critically endangered and is under the protection of the United States' Endangered Species Act. Of all seven extant Canadian beluga populations, those inhabiting eastern Hudson Bay, Ungava Bay, and the St. Lawrence River are listed as endangered.
Belugas are one of the most commonly kept cetaceans in captivity and are housed in aquariums, dolphinariums and wildlife parks in North America, Europe and Asia. They are considered charismatic because of their docile demeanour and characteristic smile, communicative nature, and supple, graceful movement.
Taxonomy
The beluga was first described in 1776 by Peter Simon Pallas. It is a member of the family Monodontidae, which is in turn part of the parvorder Odontoceti (toothed whales). The Irrawaddy dolphin was once placed in the same family, though recent genetic evidence suggests these dolphins belong to the family Delphinidae. The narwhal is the only other species within the Monodontidae. A skull has been discovered with intermediate characteristics supporting the hypothesis that hybridisation is possible between these two species.
The name of the genus, Delphinapterus, means "dolphin without fin" (from the Greek δελφίν (delphin), dolphin and απτερος (apteros), without fin) and the species name leucas means "white" (from the Greek λευκας (leukas), white). The Red List of Threatened Species gives both beluga and white whale as common names, though the former is now more popular. The English name comes from the Russian белу́га, which derives from the word белый (bélyj), meaning "white". Nowadays the word белу́га in Russian refers to the beluga sturgeon, while the whale is called almost similarly - белу́ха ("belúha").
The whale is also colloquially known as the "sea canary" on account of its high-pitched squeaks, squeals, clucks, and whistles. A Japanese researcher claimed that he taught a beluga to "talk" by using these sounds to identify three different objects, offering hope that humans may one day be able to communicate effectively with sea mammals. A similar observation has been made by Canadian researchers, where a beluga which died in 2007 "talked" when he was still a subadult. Another example is NOC, a beluga whale that could mimic the rhythm and tone of human language. Beluga whales in the wild have been reported to imitate human voices.
Evolution
Mitochondrial DNA studies have shown modern cetaceans last shared a common ancestor between 25 and 34 million years ago The superfamily Delphinoidea (which contains monodontids, dolphins and porpoises) split from other toothed whales, odontoceti, between 11 and 15 million years ago. Monodontids then split from dolphins (Delphinidae) and later from porpoises (Phocoenidae), their closest relatives in evolutionary terms. In 2017 the genome of a beluga whale was sequenced, comprising 2.327 Gbp of assembled genomic sequence that encoded 29,581 predicted genes. The authors estimated that the genome-wide sequence similarity between beluga whales and killer whales is 97.87%.
The beluga's earliest known distinctive ancestors include the prehistoric Denebola brachycephala from the late Miocene epoch (9–10 million years ago), and Bohaskaia monodontoides, from the early Pliocene (3–5 million years ago). Fossil evidence from Baja California and Virginia indicate the family once inhabited warmer waters. A fossil of the monodontid Casatia thermophila, from five million years ago, provides the strongest evidence that monodontids once inhabited warmer waters, as the fossil was found alongside fossils of tropical species such as bull and tiger sharks.
The fossil record also indicates that, in comparatively recent times, the beluga's range varied with that of the polar ice packs expanding during ice ages and contracting when the ice retreated. Counter-evidence to this theory comes from the finding in 1849 of fossilised beluga bones in Vermont in the United States, from the Atlantic Ocean. The bones were discovered during construction of the first railroad between Rutland and Burlington in Vermont, when workers unearthed the bones of a mysterious animal in Charlotte. Buried nearly below the surface in a thick blue clay, these bones were unlike those of any animal previously discovered in Vermont. Experts identified the bones as those of a beluga. Because Charlotte is over from the nearest ocean, early naturalists were at a loss to explain the presence of the bones of a marine mammal buried beneath the fields of rural Vermont.
The remains were found to be preserved in the sediments of the Champlain Sea, an extension of the Atlantic Ocean within the continent resulting from the rise in sea level at the end of the ice ages some 12,000 years ago. Today, the Charlotte whale is the official Vermont State Fossil (making Vermont the only state whose official fossil is that of a still extant animal).
Description
Its body is round, particularly when well fed, and tapers less smoothly to the head than the tail. The sudden tapering to the base of its neck gives it the appearance of shoulders, unique among cetaceans. The tail-fin grows and becomes increasingly and ornately curved as the animal ages. The flippers are broad and short—making them almost square-shaped.
Longevity
Preliminary investigations suggested a beluga's life expectancy was rarely more than 30 years. The method used to calculate the age of a beluga is based on counting the layers of dentine and dental cement in a specimen's teeth, which were originally thought to be deposited once or twice a year. The layers can be readily identified as one layer consists of opaque dense material and the other is transparent and less dense. It is therefore possible to estimate the age of the individual by extrapolating the number of layers identified and the estimated frequency with which the deposits are laid down. A 2006 study using radiocarbon dating of the dentin layers showed the deposit of this material occurs with a lesser frequency (once per year) than was previously thought. The study therefore estimated belugas can live for 70 or 80 years. However, recent studies suggest that it is unclear as to whether belugas receive a different number of layers per year depending on the age of the animal (for example young belugas may only receive an additional one layer per year), or simply just one layer per year or every other year.
Size
The species presents a moderate degree of sexual dimorphism, as the males are 25% longer than the females and are sturdier. Adult male belugas can range from , while the females measure . Males weigh between , and occasionally up to while females weigh between . They rank as mid-sized species among toothed whales.
Individuals of both sexes reach their maximum size by the time they are 10 years old. The beluga's body shape is stocky and fusiform (cone-shaped with the point facing backwards), and they frequently have folds of fat, particularly along the ventral surface. Between 40% and 50% of their body weight is fat, which is a higher proportion than for cetaceans that do not inhabit the Arctic, where fat only represents 30% of body weight. The fat forms a layer that covers all of the body except the head, and it can be up to thick. It acts as insulation in waters with temperatures between 0 and 18 °C, as well as being an important reserve during periods without food.
Colour
The adult beluga is rarely mistaken for any other species, because it is completely white or whitish-grey in colour. Calves are usually born grey, and by the time they are a month old, have turned dark grey or blue grey. They then start to progressively lose their pigmentation until they attain their distinctive white colouration, at the age of seven years in females and nine in males. The white colouration of the skin is an adaptation to life in the Arctic that allows belugas to camouflage themselves in the polar ice caps as protection against their main predators, polar bears and killer whales. Unlike other cetaceans, the belugas seasonally shed their skin. During the winter, the epidermis thickens and the skin can become yellowish, mainly on the back and fins. When they migrate to the estuaries during the summer, they rub themselves on the gravel of the riverbeds to remove the cutaneous covering.
Head and neck
Like most toothed whales, the beluga has a compartment found at the centre of the forehead that contains an organ used for echolocation called a melon, which contains fatty tissue. The shape of the beluga's head is unlike that of any other cetacean, as the melon is extremely bulbous, lobed and visible as a large frontal prominence. Another distinctive characteristic it possesses is the melon is malleable; its shape is changed during the emission of sounds. The beluga is able to change the shape of its head by blowing air around its sinuses to focus the emitted sounds. This organ contains fatty acids, mainly isovaleric acid (60.1%) and long-chain branched acids (16.9%), a very different composition from its body fat, and which could play a role in its echolocation system.
Unlike many dolphins and whales, the seven vertebrae in the neck are not fused together, allowing the animal to turn its head laterally without needing to rotate its body. This gives the head a lateral manoeuvrability that allows an improved field of view and movement and helps in catching prey and evading predators in deep water. The rostrum has about eight to ten small, blunt and slightly curved teeth on each side of the jaw and a total of 36 to 40 teeth. Belugas do not use their teeth to chew, but for catching hold of their prey; they then tear them up and swallow them nearly whole.
Belugas only have a single spiracle, which is located on the top of the head behind the melon, and has a muscular covering, allowing it to be completely sealed. Under normal conditions, the spiracle is closed and an animal must contract the muscular covering to open the spiracle. A beluga's thyroid gland is larger than that of terrestrial mammals—weighing three times more than that of a horse—which helps it to maintain a greater metabolism during the summer when it lives in river estuaries. It is the marine cetacean that most frequently develops hyperplastic and neoplastic lesions of the thyroid.
Fins
The fins retain the bony vestiges of the beluga's mammalian ancestors, and are firmly bound together by connective tissue. The fins are small in relation to the size of the body, rounded and oar-shaped and slightly curled at the tips. These versatile extremities are mainly used as a rudder to control direction, to work in synchrony with the tailfin and for agile movement in shallow waters up to deep. The fins also contain a mechanism for regulating body temperature, as the arteries feeding the fin's muscles are surrounded by veins that dilate or contract to gain or lose heat. The tailfin is flat with two oar-like lobes, it does not have any bones, and is made up of hard, dense, fibrous connective tissue. The tailfin has a distinctive curvature along the lower edge. The longitudinal muscles of the back provide the ascending and descending movement of the tailfin, which has a similar thermoregulation mechanism to the pectoral fins.
Belugas have a dorsal ridge, rather than a dorsal fin. The absence of the dorsal fin is reflected in the genus name of the species—apterus the Greek word for "wingless". The evolutionary preference for a dorsal ridge rather than a fin is believed to be an adaptation to under-ice conditions, or possibly as a way of preserving heat. The crest is hard and, along with the head, can be used to open holes in ice up to thick.
Senses
The beluga has a very specialised sense of hearing and its auditory cortex is highly developed. It can hear sounds within the range of 1.2 to 120 kHz, with the greatest sensitivity between 10 and 75 kHz, where the average hearing range for humans is 0.02 to 20 kHz. The majority of sounds are most probably received by the lower jaw and transmitted towards the middle ear. In the toothed whales, the lower jawbone is broad with a cavity at its base, which projects towards the place where it joins the cranium. A fatty deposit inside this small cavity connects to the middle ear. Toothed whales also possess a small external auditory hole a few centimetres behind their eyes; each hole communicates with an external auditory conduit and an eardrum. It is not known if these organs are functional or simply vestigial.
Belugas are able to see within and outside of water, but their vision is relatively poor when compared to dolphins. Their eyes are especially adapted to seeing under water, although when they come into contact with the air, the crystalline lens and the cornea adjust to overcome the associated myopia (the range of vision under water is short). A beluga's retina has cones and rods, which also suggests they can see in low light. The presence of cone cells indicates they can see colours, although this suggestion has not been confirmed. Glands located in the medial corner of their eyes secrete an oily, gelatinous substance that lubricates the eye and helps flush out foreign bodies. This substance forms a film that protects the cornea and the conjunctiva from pathogenic organisms.
Studies on captive animals show they seek frequent physical contact with other belugas. Areas in the mouth have been found that could act as chemoreceptors for different tastes, and they can detect the presence of blood in water, which causes them to react immediately by displaying typical alarm behaviour. Like the other toothed whales, their brains lack olfactory bulbs and olfactory nerves, which suggests they do not have a sense of smell.
Behaviour
Social structure and play
These cetaceans are highly sociable and they regularly form small groups, or pods, that may contain between two and 25 individuals, with an average of 10 members. Pods tend to be unstable, meaning individuals tend to move from pod to pod. Radio tracking has even shown belugas can start out in one pod and within a few days be hundreds of miles away from that pod. Beluga whale pods can be grouped into three categories, nurseries (which consist of mother and calves), bachelors (which consist of all males) and mixed groups. Mixed groups contain animals of both sexes. Many hundreds and even thousands of individuals can be present when the pods join in river estuaries during the summer. This can represent a significant proportion of the total population and is when they are most vulnerable to being hunted.
They are cooperative animals and frequently hunt in coordinated groups. The animals in a pod are very sociable and often chase each other as if they are playing or fighting, and they often rub against each other. Often individuals will surface and dive together in a synchronized manner, in a behavior known as milling.
In captivity, they can be seen to be constantly playing, vocalising and swimming around each other. In one case, one whale blew bubbles, while the other one popped them. There have also been reports of beluga whales copying and imitating one another, similar to a game of Simon-says. There have also been reports of them displaying physical affection, via mouth to mouth contact. They also show a great deal of curiosity towards humans and frequently approach the windows in the tanks to observe them.
Belugas also show a great degree of curiosity towards humans in the wild, and frequently swim alongside boats. They also play with objects they find in the water; in the wild, they do this with wood, plants, dead fish and bubbles they have created. During the breeding season, adults have been observed carrying objects such as plants, nets, and even the skeleton of a dead reindeer on their heads and backs. Captive females have also been observed displaying this behavior, carrying items such as floats and buoys, after they have lost a calf. Experts consider this interaction with the objects to be a substitute behavior.
In captivity, mothering behavior among belugas depends on the individual. Some mothers are extremely attentive while other mothers are so blasé, that they have actually lost their calves. In aquaria, there have been cases where dominant females have stolen calves from mothers, particularly if they have lost a calf or if they are pregnant. After giving birth, dominant females will return the calf back to their mother. Additionally, male calves will temporarily leave their mothers to interact with an adult male who can serve as a role model for the calf, before they return to their mothers. Male calves are also frequently seen interacting with each other.
Swimming and diving
Belugas are slower swimmers than the other toothed whales, such as the killer whale and the common bottlenose dolphin, because they are less hydrodynamic and have limited movement of their tail-fins, which produce the greatest thrust. They frequently swim at speeds between , although they are able to maintain a speed of 22 km/h for up to 15 min. Unlike most cetaceans, they are capable of swimming backwards. Belugas swim on the surface between 5% and 10% of the time, while for the rest of the time they swim at a depth sufficient enough to cover their bodies. They do not jump out of the water like dolphins or killer whales.
These animals usually only dive to depths to , although they are capable of diving to greater depths. Individual captive animals have been recorded at depths between 400 and 647 m below sea level, while animals in the wild have been recorded as diving to a depth of more than 700 m, with the greatest recorded depth being over 900 m. A dive normally lasts 3 to 5 minutes, but can last up to over 20 minutes. In the shallower water of the estuaries, a diving session may last around two minutes; the sequence consists of five or six rapid, shallow dives followed by a deeper dive lasting up to one minute. The average number of dives per day varies between 31 and 51.
All cetaceans, including belugas, have physiological adaptations designed to conserve oxygen while they are under water. During a dive, these animals will reduce their heart rate from 100 beats a minute to between 12 and 20. Blood flow is diverted away from certain tissues and organs and towards the brain, heart and lungs, which require a constant oxygen supply. The amount of oxygen dissolved in the blood is 5.5%, which is greater than that found in land-based mammals and is similar to that of Weddell seals (a diving marine mammal). One study found a female beluga had 16.5 L of oxygen dissolved in her blood. Lastly, the beluga's muscles contain high levels of the protein myoglobin, which stores oxygen in muscle. Myoglobin concentrations in belugas are several times greater than for terrestrial mammals, which help prevent oxygen deficiency during dives.
Beluga whales often accompany bowhead whales, for curiosity and to secure polynya feasibility to breathe as bowheads are capable of breaking through ice from underwater by headbutting.
Diet
Belugas play an important role in the structure and function of marine resources in the Arctic Ocean, as they are the most abundant toothed whales in the region. They are opportunistic feeders; their feeding habits depend on their locations and the season. For example, when they are in the Beaufort Sea, they mainly eat Arctic cod (Boreogadus saida) and the stomachs of belugas caught near Greenland were found to contain rose fish (Sebastes marinus), Greenland halibut (Reinhardtius hippoglossoides) and northern shrimp (Pandalus borealis), while in Alaska their staple diet is Coho salmon (Oncorhynchus kisutch). In general, the diets of these cetaceans consist mainly of fish; apart from those previously mentioned, other fish they feed on include capelin (Mallotus villosus), smelt, sole, flounder, herring, sculpin and other types of salmon. They also consume a great quantity of invertebrates, such as shrimp, squid, crabs, clams, octopus, sea snails, bristle worms and other deep-sea species. Belugas feed mainly in winter as their blubber is thickest in later winter and early spring, and thinnest in the fall. Inuit observation has led scientists to believe that belugas do not hunt during migration, at least in Hudson Bay.
The diet of Alaskan belugas is quite diverse and varies depending on season and migratory behavior. Belugas in the Beaufort Sea mainly feed on staghorn and shorthorn sculpin, walleye pollock, Arctic cod, saffron cod and Pacific sand lance. Shrimp are the most common invertebrate eaten, with octopus, amphipods and echiurids being other sources of invertebrate prey. The most common prey species for belugas in the Eastern Chukchi Sea appears to be shrimp, echiurid worms, cephalopods and polychaetes. The largest prey item consumed by beluga whales in the Eastern Chukchi Sea seems to be saffron cod. Beluga whales in the Eastern Bering Sea feed on a variety of fish species including saffron cod, rainbow smelt, walleye pollock, Pacific salmon, Pacific herring and several species of flounder and sculpin. The primary invertebrate consumed is shrimp. The primary prey item in regard to fish species for belugas in Bristol Bay appears to be the five species of salmon, with sockeye being the most prevalent. Smelt is also another common fish family eaten by belugas in this region. Shrimp is the most prevalent invertebrate prey item. The most common prey items for belugas in Cook Inlet appear to be salmon, cod and smelt.
Animals in captivity eat 2.5% to 3.0% of their body weight per day, which equates to . Like their wild counterparts, captive belugas were found to eat less in the fall.
Foraging on the seabed typically takes place at depths between , although they can dive to depths of in search of food. Their flexible necks provide a wide range of movement while they are searching for food on the ocean floor. Some animals have been observed to suck up water and then forcefully expel it to uncover their prey hidden in the silt on the seabed. As their teeth are neither large nor sharp, belugas must use suction to bring their prey into their mouths; it also means their prey has to be consumed whole, which in turn means it cannot be too large or the belugas run the risk of it getting stuck in their throats. They also join into coordinated groups of five or more to feed on shoals of fish by steering the fish into shallow water, where the belugas then attack them. For example, in the estuary of the Amur River, where they mainly feed on salmon, groups of six or eight individuals join to surround a shoal of fish and prevent their escape. Individuals then take turns feeding on the fish.
Reproduction
Estimations of the age of sexual maturity for beluga whales vary considerably; the majority of authors estimate males reach sexual maturity when they are between nine and fifteen years old, and females reach maturity between eight and fourteen years old. The average age at which females first give birth is 8.5 years and fertility begins to decrease when they are 25, eventually undergoing menopause, and ceasing reproductive potential with no births recorded for females older than 41. There is a slight difference on the sexual maturation period between males and females. The male beluga whales take seven to nine years to become sexually mature, while the females take four to seven years.
Female belugas typically give birth to one calf every three years. Most mating occurs from February to May, but some occurs at other times of year. The beluga may have delayed implantation. Gestation has been estimated to last 12.0 to 14.5 months, but information derived from captive females suggests a period up to 475 days (15.8 months). During the mating season, the testes of belugas double in weight. Testosterone levels increase, but seems to be independent of copulation. Copulation typically takes place between 3 and 4 AM.
Calves are born over a protracted period that varies by location. In the Canadian Arctic, calves are born between March and September, while in Hudson Bay, the peak calving period is in late June, and in Cumberland Sound, most calves are born from late July to early August. Births usually take place in bays or estuaries where the water is between 10 and 15 °C. Newborns are about long, weigh about , and are grey in colour. They are able to swim alongside their mothers immediately after birth. The newborn calves nurse under water and initiate lactation a few hours after birth; thereafter, they feed at intervals around an hour. Studies of captive females have indicated their milk composition varies between individuals and with the stage of lactation; it has an average content of 28% fat, 11% protein, 60.3% water, and less than 1% residual solids. The milk contains about 92 cal per ounce.
The calves remain dependent on their mothers for nursing for the first year, when their teeth appear. After this, they start to supplement their diets with shrimp and small fish. The majority of the calves continue nursing until they are 20 months old, although occasionally lactation can continue for more than two years, and lactational anoestrus may not occur. Alloparenting (care by females different from the mother) has been observed in captive belugas, including spontaneous and long-term milk production. This suggests this behaviour, which is also seen in other mammals, may be present in belugas in the wild.
Hybrids have been documented between the beluga and the narwhal (specifically offspring conceived by a beluga father and a narwhal mother), as one, perhaps even as many as three, such hybrids were killed and harvested during a sustenance hunt. Whether or not these hybrids could breed remains unknown. The unusual dentition seen in the single remaining skull indicates the hybrid hunted on the seabed, much as walruses do, indicating feeding habits different from those of either parent species.
Communication and echolocation
Belugas use sounds and echolocation for movement, communication, to find breathing holes in the ice, and to hunt in dark or turbid waters. They produce a rapid sequence of clicks that pass through the melon, which acts as an acoustic lens to focus the sounds into a beam that is projected forward through the surrounding water. These sounds spread through the water at a speed of nearly 1.6 km per second, some four times faster than the speed of sound in air. The sound waves reflect from objects and return as echoes that are heard and interpreted by the animal. This enables them to determine the distance, speed, size, shape and the object's internal structure within the beam of sound. They use this ability when moving around thick Arctic ice sheets, to find areas of unfrozen water for breathing, or air pockets trapped under the ice.
Some evidence indicates that belugas are highly sensitive to noise produced by humans. In one study, the maximum frequencies produced by an individual located in San Diego Bay, California, were between 40 and 60 kHz. The same individual produced sounds with a maximum frequency of 100 to 120 kHz when transferred to Kaneohe Bay in Hawaii. The difference in frequencies is thought to be a response to the difference in environmental noise in the two areas. In special circumstances, beluga whale sounds have been reported to resemble human speech.
These animals communicate using sounds of high frequency; their calls can sound like bird songs, so belugas were nicknamed "canaries of the sea". Like the other toothed whales, belugas do not possess vocal cords and the sounds are probably produced by the movement of air between the nasal sacks, which are located near to the blowhole.
Belugas are among the most vocal cetaceans. They use their vocalisations for echolocation, during mating and for communication. They possess a large repertoire, emitting up to 11 different sounds, such as cackles, whistles, trills and squawks. They make sounds by grinding their teeth or splashing, but they rarely use body language.
There is debate as to whether cetacean vocalizations can constitute a language. A study conducted in 2015 determined that European beluga signals share physical features comparable to vowels. These sounds were found to be stable throughout time, but varied among different geographical locations. The further away the populations were from each other, the more varied the sounds were in relation to one another.
Distribution
The beluga inhabits a discontinuous circumpolar distribution in Arctic and sub-Arctic waters. During the summer, they can mainly be found in deep waters ranging from 76°N to 80°N, particularly along the coasts of Alaska, northern Canada, western Greenland and northern Russia. The southernmost extent of their range includes isolated populations in the St. Lawrence River in the Atlantic, and the Amur River delta, the Shantar Islands and the waters surrounding Sakhalin Island in the Sea of Okhotsk.
Migration
Belugas have a seasonal migratory pattern. Migration patterns are passed from parents to offspring. Some travel as far as per year. When the summer sites become blocked with ice during the autumn, they move to spend the winter in the open sea alongside the pack ice or in areas covered with ice, surviving by using polynyas to surface and breathe. In summer after the sheet ice has melted, they move to coastal areas with shallower water ( deep), although sometimes they migrate towards deeper waters (deeper than ). In the summer, they occupy estuaries and the waters of the continental shelf, and, on occasion, they even swim up the rivers. A number of incidents have been reported where groups or individuals have been found hundreds or even thousands of kilometres from the ocean. One such example comes from June 9, 2006, when a young beluga carcass was found in the Tanana River near Fairbanks in central Alaska, nearly from the nearest ocean habitat. Belugas sometimes follow migrating fish, leading Alaska state biologist Tom Seaton to speculate it had followed migrating salmon up the river at some point in the previous autumn. The rivers they most often travel up include: the Northern Dvina, the Mezen, the Pechora, the Ob and the Yenisei in Asia; the Yukon and the Kuskokwim in Alaska, and the Saint Lawrence in Canada. Spending time in a river has been shown to stimulate an animal's metabolism and facilitates the seasonal renewal of the epidermal layer. In addition, the rivers represent a safe haven for newborn calves where they will not be preyed upon by killer whales. Calves often return to the same estuary as their mother in the summer, meeting her sometimes even after becoming fully mature. However, not all beluga whale populations summer in estuaries. Belugas from the Beaufort Sea stock were found to summer along the Eastern Beaufort Sea shelf, Amundsen Gulf and slope regions north and west of Banks Island, in addition to core areas in the Mackenzie River Estuary. Male belugas have been observed summering in deeper waters along Viscount Melville Sound, in depths of up to . The bulk of Eastern Chukchi Sea belugas summer over Barrow canyon.
The migration season is relatively predictable, as it is basically determined by the amount of daylight and not by other variable physical or biological factors, such as the condition of the sea ice. Vagrants may travel further south to areas such as Irish and Scottish waters, the islands of Orkney and Hebrides, and to Japanese waters. There had been several vagrant individuals that have demonstrated seasonal residencies at Volcano Bay, and a unique whale were used to return annually to areas adjacent to Shibetsu in Nemuro Strait in the 2000s. On rarer occasions, individuals of vagrancy can reach the Korean Peninsula. A few other individuals have been confirmed to return to the coasts of Hokkaido, and one particular individual became a resident in brackish waters of Lake Notoro since in 2014.
Some populations are not migratory and certain resident groups will stay in well-defined areas, such as in Cook Inlet, the estuary of the Saint Lawrence River and Cumberland Sound. The population in Cook Inlet stays in the waters furthest inside the inlet during the summer until the end of autumn. Then during the winter, they disperse to the deeper water in the center of the inlet, but without completely leaving it.
In April, the animals that spend the winter in the center and southwest of the Bering Sea move to the north coast of Alaska and the east coast of Russia. The populations living in the Ungava Bay and the eastern and western sides of Hudson Bay overwinter together beneath the sea ice in Hudson Strait. Whales in James Bay that spend winter months within the basin, could be a distinct group from those in Hudson Bay. The populations of the White Sea, the Kara Sea and the Laptev Sea overwinter in the Barents Sea. In the spring, the groups separate and migrate to their respective summer sites.
Habitat
Belugas exploit a varied range of habitats; they are most commonly seen in shallow waters close to the coast, but they have also been reported to live for extended periods in deeper water, where they feed and give birth to their young.
In coastal areas, they can be found in coves, fjords, canals, bays and shallow waters in the Arctic Ocean that are continuously lit by sunlight. They are also often seen during the summer in river estuaries, where they feed, socialize and give birth to young. These waters usually have a temperature between 8 and 10 °C. The mudflats of Cook Inlet in Alaska are a popular location for these animals to spend the first few months of summer. In the eastern Beaufort Sea, female belugas with their young and immature males prefer the open waters close to land, while the adult males live in waters covered by ice near the Canadian Arctic Archipelago. The younger males and females with slightly older young can be found nearer to the ice shelf. Generally, the use of different habitats in summer reflects differences in feeding habits, risk from predators and reproductive factors for each of the subpopulations.
Population
There are currently 22 stocks of beluga whales recognized:
James Bay – 14,500 individuals (belugas remain here all year round)
Western Hudson Bay – 55,000 individuals
Eastern Hudson Bay – 3,400–3,800 individuals
Cumberland Sound – 1,151 individuals
Ungava Bay – 32 individuals (maybe functionally extinct)
St. Lawrence River Estuary – 889 individuals
Eastern Canadian Arctic – 21,400 individuals
Southwest Greenland – Extinct
Eastern Chukchi Sea – 20,700 individuals
Eastern Bering Sea – 7,000–9,200 individuals
Eastern Beaufort Sea – 39,300 individuals
Bristol Bay – 2,000–3,000 individuals
Cook Inlet – 300 individuals
White Sea – 5,600 individuals
Kara Sea/Laptev Sea/Barents Sea – Data Deficient
Ulbansky – 2,300
Anadyr – 3,000
Shelikhov – 2,666
Sakhalin/Amur – 4,000 individuals
Tugurskiy – 1,500 individuals
Udskaya – 2,500 individuals
Svalbard – 549 individuals
The Yakutat Bay belugas are not considered to be a true stock because they have only been present in these waters since the 1980s, and are believed to be of Cook Inlet origin. It is estimated that less than 20 whales inhabit the bay year-round. Overall the beluga population is estimated to be 150,000–200,000 animals.
Threats
Hunting
The native populations of the Arctic in Alaska, Canada, Greenland and Russia hunt belugas, for both consumption and profit. Belugas have been easy prey for hunters due to their predictable migration patterns and the high population density in estuaries and surrounding coastal areas during the summer.
Present
The number of animals killed is about 1,000 per year, (see table below. and its sources). Beluga whale hunting quotas in Canada and the United States are established using the Potential Biological Removal equation PBR = Nmin * 0.5 * Rmax * FR, to determine what constitutes a sustainable hunt. Nmin represents a conservative estimation of the population size, Rmax, represents the maximum rate of population increase and FR represents the recovery factor.
Hunters in Hudson's Bay rarely eat the meat. They give a little to dogs, and leave the rest for wild animals. Other areas may dry the meat for later consumption by humans. In Greenland the skin (muktuk) is sold commercially to fish factories, and in Canada to other communities. An average of one or two vertebrae and one or two teeth per beluga are carved and sold. One estimate of the annual gross value received from Beluga hunts in Hudson Bay in 2013 was for 190 belugas, or per beluga. However, the net income, after subtracting costs in time and equipment, was a loss of per person. Hunts receive subsidies, but they continue as a tradition, rather than for the money, and the economic analysis noted that whale watching may be an alternate revenue source. Of the gross income, was for skin and meat, to replace beef, pork and chickens which would otherwise be bought. was received for carved vertebrae and teeth.
Russia now harvests 5 to 30 belugas per year for meat and captures an additional 20 to 30 per year for live export to Chinese aquaria. However, in 2018, 100 were illegally captured for live export.
Previous levels of commercial whaling have put the species in danger of extinction in areas such as Cook Inlet, Ungava Bay, the St. Lawrence River and western Greenland. Continued hunting by the native peoples may mean some populations will continue to decline. Northern Canadian sites are the focus of discussions between local communities and the Canadian government, with the objective of permitting sustainable hunting that does not put the species at risk of extinction.
The total amount of landed (defined as belugas successfully hunted and retrieved) belugas averages 275 in regard to the Bering, Chukchi and Beaufort stocks from 1987 to 2006. The average annual landed harvest of belugas in the Beaufort Sea consisted of 39 individuals while the Chukchi harvest averaged 62 individuals. Bristol bay's annual average landed harvest was 17 while the Bering Sea's was 152. Statistical studies have demonstrated that subsistence hunting in Alaska did not significantly impact the population of the Alaskan beluga whale stocks. The number of belugas struck and lost did not seem to profoundly impact Chukchi and Bering Sea belugas.
Past
Commercial whaling by European, American and Russian whalers during the 18th and 19th centuries decreased beluga populations in the Arctic. The animals were hunted for their meat and blubber, while the Europeans used the oil from the melon as a lubricant for clocks, machinery and lighting in lighthouses. Mineral oil replaced whale oil in the 1860s, but into the early 20th century the cured skin was still used to make horse harnesses and machine belts for saw mills and shoelaces. The cured skin is the only cetacean skin that is sufficiently thick to be used as leather, and was used to manufacture some of the first bulletproof vests.
Russia had large hunts, peaking in the 1930s at 4,000 per year and the 1960s at 7,000 per year, for a total of 86,000 from 1915 to 2014. Canada hunted a total of 54,000 from 1731 to 1970. Between 1868 and 1911, Scottish and American whalers killed more than 20,000 belugas in Lancaster Sound and Davis Strait.
During the 1920s, fishermen in the Saint Lawrence River estuary considered belugas to be a threat to the fishing industry, as they eat large quantities of cod, salmon, tuna and other fish caught by the local fishermen. The presence of belugas in the estuary was, therefore, considered to be undesirable; in 1928, the Government of Quebec offered a reward of 15 dollars for each dead beluga. The Quebec Department of Fisheries launched a study into the influence of these cetaceans on local fish populations in 1938. The unrestricted killing of belugas continued into the 1950s, when the supposed voracity of the belugas was found to be overestimated and did not adversely affect fish populations. L'Isle-aux-Coudres is the setting for the classic 1963 National Film Board of Canada documentary Pour la suite du monde, which depicts a one-off resurrection of the beluga hunt; one animal is caught live, and transported by truck to an aquarium in the big city. The method of capture is akin to dolphin drive hunting.
Beluga catches by location
Predation
During the winter, belugas commonly become trapped in the ice without being able to escape to open water, which may be several kilometres away. Polar bears take particular advantage of these situations and are able to locate the belugas using their sense of smell. The bears swipe at the belugas and drag them onto the ice to eat them. They are able to capture large individuals in this way; in one documented incident, a bear weighing between 150 and 180 kg was able to capture a beluga that weighed 935 kg.
Killer whales hunt and eat both young and adult belugas. They live in all the seas of the world and share the same habitat as belugas in the sub-Arctic region. Attacks on belugas by killer whales have been reported in the waters of Greenland, Russia, Canada and Alaska. A number of killings have been recorded in Cook Inlet, and experts are concerned the predation by killer whales will impede the recovery of this sub-population, which has already been badly depleted by hunting. The killer whales arrive at the beginning of August, but the belugas are occasionally able to hear their presence and evade them. The groups near to or under the sea ice have a degree of protection, as the killer whale's large dorsal fin, up to 2 m in length, impedes their movement under the ice and does not allow them to get sufficiently close to the breathing holes in the ice. Beluga whale behavior under killer whale predation makes them vulnerable to hunters. When killer whales are present, large numbers of beluga whales congregate in the shallows for protection, which allows them to be hunted in droves.
Contamination
The beluga is considered an excellent sentinel species (indicator of environment health and changes), because it is long-lived, at the top of the food web, bears large amounts of fat and blubber, relatively well-studied for a cetacean, and still somewhat common.
Human pollution can be a threat to belugas' health when they congregate in river estuaries. Chemical substances such as DDT and heavy metals such as lead, mercury and cadmium have been found in individuals of the Saint Lawrence River population. Local beluga carcasses contain so many contaminants, they are treated as toxic waste. Levels of polychlorinated biphenyls between 240 and 800 ppm have been found in belugas' brains, liver and muscles, with the highest levels found in males. These levels are significantly greater than those found in Arctic populations. These substances have a proven adverse effect on these cetaceans, as they cause cancers, reproductive diseases and the deterioration of the immune system, making individuals more susceptible to pneumonias, ulcers, cysts, tumours and bacterial infections. Although the populations that inhabit the river estuaries run the greatest risk of contamination, high levels of zinc, cadmium, mercury and selenium have also been found in the muscles, livers and kidneys of animals that live in the open sea. The concentration of mercury in Beaufort Sea belugas tripled from the 1980s to the 1990s, but has decreased in Beaufort belugas as of the 21st century, possibly due to changes in dietary preference. Larger body sized belugas tend to have more mercury than smaller sized belugas, because they spend more time offshore, hunting prey such as cod and shrimp, which have more mercury.
From a sample of 129 beluga adults from the Saint Lawrence River examined between 1983 and 1999, a total of 27% had suffered cancer. This is a higher percentage than that documented for other populations of this species and is much higher than for other cetaceans and for the majority of terrestrial mammals; in fact, the rate is only comparable to the levels found in humans and some domesticated animals. For example, the rate of intestinal cancer in the sample is much higher than for humans. This condition is thought to be directly related to environmental contamination, in this case by polycyclic aromatic hydrocarbons, and coincides with the high incidence of this disease in humans residing in the area. The prevalence of tumours suggests the contaminants identified in the animals that inhabit the estuary are having a direct carcinogenic effect or they are at least causing an immunological deterioration that is reducing the inhabitants' resistance to the disease.
Indirect human disturbance may also be a threat. While some populations tolerate small boats, most actively try to avoid ships. Whale-watching has become a booming activity in the St. Lawrence and Churchill River areas, and acoustic contamination from this activity appears to have an effect on belugas. For example, a correlation appears to exist between the passage of belugas across the mouth of the Saguenay River, which has decreased by 60%, and the increase in the use of recreational motorboats in the area. A dramatic decrease has also been recorded in the number of calls between animals (decreasing from 3.4 to 10.5 calls/min to 0 or <1) after exposure to the noise produced by ships, the effect being most persistent and pronounced with larger ships such as ferries than with smaller boats. Belugas can detect the presence of large ships (for example icebreakers) up to 50 km away, and they move rapidly in the opposite direction or perpendicular to the ship following the edge of the sea ice for distances of up to 80 km to avoid them. The presence of shipping produces avoidance behaviour, causing deeper dives for feeding, the break-up of groups, and asynchrony in dives.
Pathogens
As with any animal population, a number of pathogens cause death and disease in belugas, including viruses, bacteria, protozoans and fungi, which mainly cause skin, intestinal and respiratory infections.
Papillomaviruses, herpesviruses, and encephalitis caused by the protozoan Sarcocystis have been found in belugas in the Saint Lawrence River. Cases have been recorded of ciliate protozoa colonising the spiracle of certain individuals, but they are not thought to be pathogens or are not very harmful. The bacterium Erysipelothrix rhusiopathiae, which probably comes from eating infected fish, poses a threat to belugas kept in captivity, causing anorexia and dermal plaques and lesions that can lead to sepsis. This condition can cause death if it is not diagnosed and treated in time with antibiotics such as ciprofloxacin. A study of infections caused by parasitic worms in a number of individuals of both sexes found the presence of larvae from a species from the genus Contracaecum in their stomachs and intestines, Anisakis simplex in their stomachs, Pharurus pallasii in their ear canals, Hadwenius seymouri in their intestines and Leucasiella arctica in their rectums.
Relationship with humans
Captivity
Belugas were among the first whale species to be kept in captivity. The first beluga was shown at Barnum's Museum in New York City in 1861. For most of the 20th century, Canada was the predominant source for belugas destined for exhibition. Throughout the early 1960s, belugas were taken from the St. Lawrence River estuary. In 1967, the Churchill River estuary became the main source from which belugas were captured. This continued until 1992, when the practice was banned. Since then, Russia has become the largest provider. Individuals are caught in the Amur River delta and the far eastern seas of the country, and then are either transported domestically to aquaria in Moscow, St. Petersburg and Sochi, or exported to foreign nations, including China and formerly Canada. Canada has now banned the practice of holding new animals in captivity.
To provide some enrichment while in captivity, aquaria train belugas to perform behaviours for the public and for medical exams, such as blood draws, ultrasound, providing toys, and allowing the public to play recorded or live music.
Between 1960 and 1992, the United States Navy carried out a program that included the study of marine mammals' abilities with echolocation, with the objective of improving the detection of underwater objects. The program started with dolphins, but a large number of belugas were also used from 1975 onwards. The program included training them to carry equipment and material to divers working under water, the location of lost objects, surveillance of ships and submarines, and underwater monitoring using cameras held in their mouths. A similar program was implemented by the Soviet Navy during the Cold War, in which belugas were also trained for antimining operations in Arctic waters. It is possible this program continues within the Russian Navy, as on April 24, 2019, a tame beluga whale wearing a Russian equipment harness was found by fishermen near the Norwegian island of Ingøya.
Belugas released from captivity have difficulties adapting to life in the wild, but if not fed by humans they may have a chance to join a group of wild belugas and learn to feed themselves, according to Audun Rikardsen of the University of Tromsø.
In 2019, a sanctuary in Iceland was established for two belugas, Little White and Little Grey, that retired from a marine park in China. The Sea Life Trust Beluga Whale Sanctuary was created with support from Merlin Entertainments and Whale and Dolphin Conservation (WDC). Merlin bought the park in 2012, as part of an Australian chain, and it is one of their largest aquaria. Merlin has a policy against captive cetaceans, so they sponsored a 32,000-square-metre sea pen as a sanctuary. The 12-year-old belugas, caught in Russia and raised in captivity, do not know how to live in the wild. The cost is variously listed as ISK 3,000,000 (US$24,000) or US$27,000,000. Merlin was owned until 2015 by Blackstone Group, which also owned SeaWorld until selling its last stake in 2017 to a Chinese company which will use SeaWorld's expertise to expand in China; SeaWorld still keeps belugas in captivity.
Belugas are displayed across North America, Europe and Asia. As of 2006, 58 belugas were held in captivity in Canada and the United States, and 42 deaths in US captivity had been reported up to that time. A single specimen costs up to US$100,000, although the price has now dropped to US$70,000. As of January 2018, according to the nonprofit Ceta Base, which tracks belugas and dolphins under human care, there were 81 captive belugas in Canada and the United States, and unknown numbers in the rest of the world. The beluga's popularity with visitors reflects its attractive colour and its range of facial expressions. The latter is possible because while most cetacean "smiles" are fixed, the extra movement afforded by the beluga's unfused cervical vertebrae allows a greater range of apparent expression.
Most belugas found in aquaria are caught in the wild, as captive-breeding programs have not had much success so far. For example, despite best efforts, as of 2010, only two male whales had been successfully used as stud animals in the Association of Zoos and Aquariums beluga population, Nanuq at SeaWorld San Diego and Naluark at the Shedd Aquarium in Chicago, US. Nanuq has fathered 10 calves, five of which survived birth. Naluark at Shedd Aquarium has fathered four living offspring. Naluark was relocated to the Mystic Aquarium in the hope that he would breed with two of their females, but he did not, and in 2016 he was moved to SeaWorld Orlando. The first beluga calf born in captivity in Europe was born in L'Oceanogràfic marine park in Valencia, Spain, in November 2006. However, the calf died 25 days later after suffering metabolic complications, infections and not being able to feed properly. A second calf was born on 16 November 2016, and was successfully maintained by artificial feeding based on enriched milk.
In 2009 during a free-diving competition in a tank of icy water in Harbin, China, a captive beluga brought a cramp-paralysed diver from the bottom of the pool up to the surface by holding her foot in its mouth, saving the diver's life.
Films which have publicised issues of beluga welfare include Born to Be Free, Sonic Sea, and Vancouver Aquarium Uncovered.
Whale watching
Whale watching has become an important activity in the recovery of the economies of towns in Quebec and Hudson Bay, near the Saint Lawrence and Churchill Rivers. The best time to see belugas is during the summer, when they meet in large numbers in the estuaries of the rivers and in their summer habitats. The animals are easily seen due to their high numbers and their curiosity regarding the presence of humans.
However, the boats' presence poses a threat, as it distracts them from important activities such as feeding, social interaction and reproduction. In addition, the noise produced by the motors has an adverse effect on their auditory function and reduces their ability to detect their prey, communicate and navigate. To protect these marine animals during whale-watching activities, the US National Oceanic and Atmospheric Administration has published a "Guide for observing marine life". The guide recommends boats carrying the whale watchers keep their distance from the cetaceans and it expressly prohibits chasing, harassing, obstructing, touching, or feeding them.
Some regular migrations do occur into Russian EEZ of Sea of Japan such as to Rudnaya Bay, where diving with wild belugas became a less-known but popular attraction.
On 25 September 2018, a beluga was sighted in the Thames Estuary and near towns along the Kent side of the Thames, being nicknamed Benny by newspapers. The whale, who was noticed by conservationists to be traveling alone, appeared to be separated from the rest of its group, and is thought to be a lost individual. Subsequent sightings were reported on the following day, and continued into 2019, when local experts concluded that Benny had left the estuary.
On 13 May 2021, two beluga whales were sighted in waters around Prince Edward Island, Atlantic Canada. One whale entered the Charlottetown Harbour and travelled up the Hillsborough River to Mount Stewart, Prince Edward Island. As of 30 May the whale was still sighted in the area.
In August 2022, a beluga was found on the river Seine, France.
Human speech
Male belugas in captivity can mimic the pattern of human speech, several octaves lower than typical whale calls. One captive male beluga named NOC, pronounced "No-see", caused a diver in a tank with him to surface by imitating orders to get out of the water. Subsequent recordings confirmed that NOC had become skilled at imitating the patterns and frequency of human speech, and intentionally altered his normal methods of vocalization to achieve these sounds. After several years, he ceased making these sounds.
Conservation status
Prior to 2008, the beluga was listed as "vulnerable" by the International Union for Conservation of Nature (IUCN), a higher level of concern. The IUCN cited the stability of the largest sub-populations and improved census methods that indicate a larger population than previously estimated. In 2008, the beluga was reclassified as "near threatened" by the IUCN due to uncertainty about threats to their numbers and the number of belugas over parts of its range (especially the Russian Arctic), and the expectation that if current conservation efforts cease, especially hunting management, the beluga population is likely to qualify for "threatened" status within five years. In June 2017, its status was reassessed to "least concern".
There are about 21 sub-populations of beluga whales and it is estimated that 200,000 individuals still exist, which are listed as Least Concern on the IUCN Red List. However, the nonmigratory Cook Inlet sub-population off the Gulf of Alaska is a separate sub-population that is listed as "critically endangered" by the IUCN as of 2006 and as "endangered" under the Endangered Species Act as of October 2008. This was primarily due to unregulated overharvesting of beluga whales prior to 1998. The population has remained relatively consistent, though the reported harvest has been small. As of 2016, the estimated abundance of the endangered Cook Inlet population was 293 individuals. The most recent estimate in 2018 by NOAA Fisheries suggested that the population declined to 279 individuals.
Despite beluga whales not being threatened overall, sub-populations are being listed as critically endangered and are facing increased mortality from human actions. For example, even though commercial hunting is now banned due to the Marine Mammal Protection Act, beluga whales are still being hunted to preserve the livelihood of native Alaskan communities. The IUCN and NOAA Fisheries cite habitat degradation, oil and gas drilling, underwater noise, harvesting for consumption and climate change as threats to the prolonged survival of beluga whale sub-populations.
Beluga whale populations are currently being harvested at levels which are not sustainable and it is difficult for those harvesting beluga whales to know which sub-population they are from. Because there is little protection of sub-populations, harvest will need to be managed to ensure sub-populations will survive long into the future to discover the importance of their migratory patterns and habitat use.
Beluga whales, like most other arctic species, are being faced with alteration of their habitat due to climate change and melting arctic ice. Changes in sea-ice has resulted in changes in the area used by Chukchi belugas, since belugas spent less time in close proximity to the ice edge in comparison to previous years. Additionally, Chukchi Sea belugas spent a prolonged amount of time in Barrow Canyon on the Beaufort Sea side in October. Chukchi sea belugas also appear to be spending more time in deeper water presently, as opposed to the 1990s. Belugas also seemed to be taking longer and deeper dives. A hypothesis as to why this might be the case is an up-welling of rich Atlantic water in the Beaufort Sea may result in concentrated prey items like Arctic cod. The fall migration of Chukchi belugas is later, although summer and fall habitat selection has not changed. Fall migration of Chukchi belugas appears to be correlated with Beaufort Sea freeze up.
It is hypothesized that beluga whales utilize ice as protection from killer whale predation or for feeding on schools of fish. Killer whales can penetrate further into the Arctic and remain in arctic waters for a longer period of time due to reductions in sea ice. For example, residents in Kotzebue, have reported that killer whales have been sighted more frequently in Kotzebue Sound.
As annual ice cover declines, humans may gain access and disrupt beluga whale habitats. For example, the number of vessels in the Arctic for gas and oil exploration, fishing, and commercial shipping has already increased and a continuous trend may lead to higher risks of injuries and deaths for beluga whales.
In addition, it is possible that beluga whales may face by an increased risk of entrapment from leads and cracks freezing, due to the erratic nature of climate change. Abrupt changes in weather can cause these leads and cracks to freeze ultimately causing the whales to die of suffocation. An increase in urbanization will likely lead to higher concentrations of toxic pollutants in the blubber of beluga whales since they are at the top of the food chain and are affected by bio-accumulation. Loss of sea ice and a change in ocean temperatures may also affect the distribution and composition of prey or affect their competition. There is also some evidence that climate change can affect males and females differently. Since 1983, belugas have been increasingly scarce in Kotzebue sound. However, in 2007, several hundred whales were spotted in the sound, with over 90% of the whales being male. However, more research needs to be conducted to understand how climate change affects beluga whale sex aggregation.
Legal protection
The US Congress passed the Marine Mammal Protection Act of 1972, outlawing the persecution and hunting of all marine mammals within US coastal waters. The act has been amended a number of times to permit subsistence hunting by native peoples, temporary capture of restricted numbers for research, education and public display, and to decriminalise the accidental capture of individuals during fishing operations. The act also states that all whales in US territorial waters are under the jurisdiction of the National Marine Fisheries Service, a division of NOAA.
To prevent hunting, belugas are protected under the 1986 International Moratorium on Commercial Whaling; however, hunting of small numbers of belugas is still allowed. Since it is very difficult to know the exact population of belugas because their habitats include inland waters away from the ocean, they easily come in contact with oil and gas development centres. To prevent whales from coming in contact with industrial waste, the Alaskan and Canadian governments are relocating sites where whales and waste come in contact.
The beluga whale is listed on appendix II of the Convention on the Conservation of Migratory Species of Wild Animals (CMS). It is listed on appendix II as it has an unfavourable conservation status or would benefit significantly from international co-operation organised by tailored agreements. All toothed whales are protected under the CITES that was signed in 1973 to regulate the international import and export of certain species.
The isolated beluga population in the Saint Lawrence River has been legally protected since 1983. In 1988 Canadian Department of Fisheries and Oceans and Environment Canada, a governmental agency that supervises national parks, implemented the Saint Lawrence Action Plan with the aim of reducing industrial contamination by 90% by 1993; as of 1992, the emissions had been reduced by 59%. The population of the St. Lawrence belugas decreased from 10,000 in 1885 to around 1,000 in the 1980 and around 900 in 2012.
Conservation research in managed care facilities
As of 2015, there were 33 individuals housed in managed care facilities in North America. These facilities are members of the Association of Zoos and Aquariums, aiming to understand the complex reproductive physiology of this species to improve their conservation. With the extreme difficulty of studying beluga whales in the wild and the lack of ability to collect biological samples or perform examinations on individuals, managed care facilities play a critical role.
Managed care facilities in North America have been able to work cooperatively to build upon the research of beluga whale reproduction and have made remarkable advances. Using operant conditioning, these facilities have trained beluga whales for voluntary biological sampling and examinations. Blood, urine, and blow samples have all been collected for longitudinal hormone monitoring studies.
In addition, beluga whales have undergone semen collection, body temperature data collection, reproductive tract examinations via transabdominal ultrasound, and endoscopic exams. With new technology, the reproductive characteristics of both the female and male beluga whale have been accurately described and has benefited captive breeding programs globally.
As more research is done, the management of beluga whales in managed care facilities can be greatly improved and may even help develop other cetacean breeding and contraceptive programs, such as that of the bottlenose dolphin. Through fetal health and gestation monitoring, facilities can be more equipped to deal with pregnant animals as well. While training has been done to collect beluga whale semen, only few facilities have been able to successfully do so as both saltwater and urine contamination need to be avoided. Improvement of this process will help increase the success of captive breeding programs.
Cultural references
Pour la suite du monde, is a Canadian documentary film released in 1963 about traditional beluga hunting carried out by the inhabitants of L'Isle-aux-Coudres on the Saint Lawrence River.
The children's singer Raffi released an album called Baby Beluga in 1980. The album starts with the sound of whales communicating, and includes songs representing the ocean and whales playing. The song "Baby Beluga" was composed after Raffi saw a recently born beluga calf in Vancouver Aquarium.
The fuselage design of the Airbus Beluga, one of the world's biggest cargo planes, is very similar to that of a beluga. It was originally called the Super Transporter, but the nickname Beluga became more popular and was then officially adopted. The company paints the 2019 Beluga XL version to emphasize the plane's similarity to the Beluga whale.
In the 2016 Disney/Pixar animated film Finding Dory, the sequel to Finding Nemo (2003), the character Bailey is a beluga whale and its echolocation abilities are a significant part of the plot.
| Biology and health sciences | Toothed whale | Animals |
231030 | https://en.wikipedia.org/wiki/Baleen%20whale | Baleen whale | Baleen whales (), also known as whalebone whales, are marine mammals of the parvorder Mysticeti in the infraorder Cetacea (whales, dolphins and porpoises), which use keratinaceous baleen plates (or "whalebone") in their mouths to sieve planktonic creatures from the water. Mysticeti comprises the families Balaenidae (right and bowhead whales), Balaenopteridae (rorquals), Eschrichtiidae (the gray whale) and Cetotheriidae (the pygmy right whale). There are currently 16 species of baleen whales. While cetaceans were historically thought to have descended from mesonychians, molecular evidence instead supports them as a clade of even-toed ungulates (Artiodactyla). Baleen whales split from toothed whales (Odontoceti) around 34 million years ago.
Baleen whales range in size from the and pygmy right whale to the and blue whale, the largest known animal to have ever existed. They are sexually dimorphic. Baleen whales can have streamlined or large bodies, depending on the feeding behavior, and two limbs that are modified into flippers. The fin whale is the fastest baleen whale, recorded swimming at . Baleen whales use their baleen plates to filter out food from the water by either lunge-feeding or skim-feeding. Baleen whales have fused neck vertebrae, and are unable to turn their heads at all. Baleen whales have two blowholes. Some species are well adapted for diving to great depths. They have a layer of fat, or blubber, under the skin to keep warm in the cold water.
Although baleen whales are widespread, most species prefer the colder waters of the Arctic and Antarctic. Gray whales are specialized for feeding on bottom-dwelling crustaceans. Rorquals are specialized at lunge-feeding, and have a streamlined body to reduce drag while accelerating. Right whales skim-feed, meaning they use their enlarged head to effectively take in a large amount of water and sieve the slow-moving prey. Males typically mate with more than one female (polygyny), although the degree of polygyny varies with the species. Male strategies for reproductive success vary between performing ritual displays (whale song) or lek mating. Calves are typically born in the winter and spring months and females bear all the responsibility for raising them. Mothers fast for a relatively long period of time over the period of migration, which varies between species. Baleen whales produce a number of infrasonic vocalizations, notably the songs of the humpback whale.
The meat, blubber, baleen, and oil of baleen whales have traditionally been used by the indigenous peoples of the Arctic. Once relentlessly hunted by commercial industries for these products, cetaceans are now protected by international law. These protections have allowed their numbers to recover. However, the North Atlantic right whale is ranked critically endangered by the International Union for Conservation of Nature. Besides hunting, baleen whales also face threats from marine pollution and ocean acidification. It has been speculated that man-made sonar results in strandings. They have rarely been kept in captivity, and this has only been attempted with juveniles or members of one of the smallest species.
Taxonomy
Baleen whales are cetaceans classified under the parvorder Mysticeti, and consist of four extant families: Balaenidae (right whales), Balaenopteridae (rorquals), Eschrichtiidae (gray whale) and Cetotheriidae (pygmy right whale). Balaenids are distinguished by their enlarged head and thick blubber, while rorquals and gray whales generally have a flat head, long throat pleats, and are more streamlined than Balaenids. Rorquals also tend to be longer than the latter. Cetaceans (whales, dolphins, and porpoises) and artiodactyls are now classified under the order Cetartiodactyla, often still referred to as Artiodactyla (given that the cetaceans are deeply nested with the artiodactyls). The closest living relatives to baleen whales are toothed whales both from the infraorder Cetacea.
Classification
Balaenidae consists of two genera: Eubalaena (right whales) and Balaena (the bowhead whale, B. mysticetus). Balaenidae was thought to have consisted of only one genus until studies done through the early 2000s reported that bowhead whales and right whales are morphologically (different skull shape) and phylogenically different. According to a study done by H. C. Rosenbaum (of the American Museum of Natural History) and colleagues, the North Pacific (E. japonica) and Southern right (E. australis) whales are more closely related to each other than to the North Atlantic right whale (E. glacialis).
Cetotheriidae consists of only one living member: the pygmy right whale (Caperea marginata). The first descriptions date back to the 1840s of bones and baleen plates resembling a smaller version of the right whale, and was named Balaena marginata. In 1864, it was moved into the genus Caperea after a skull of another specimen was discovered. Six years later, the pygmy right whale was classified under the family Neobalaenidae. Despite its name, the pygmy right whale is more genetically similar to rorquals and gray whales than to right whales. A study published in 2012, based on bone structure, moved the pygmy right whale from the family Neobalaenidae to the family Cetotheriidae, making it a living fossil; Neobalaenidae was demoted to subfamily level as Neobalaeninae.
Rorquals consist of three genera (Balaenoptera, Megaptera, and Eschrichtius) and 11 species: the fin whale (B. physalus), the Sei whale (B. borealis), Bryde's whale (B. brydei), Eden's whale (B. edeni), Rice's whale (B. ricei), the blue whale (B. musculus), the common minke whale (B. acutorostrata), the Antarctic minke whale (B. bonaerensis), Omura's whale (B. omurai), the humpback whale (M. novaeangliae), and the gray whale (E. robustus). In a 2012 review of cetacean taxonomy, Alexandre Hassanin (of the Muséum National d'Histoire Naturelle) and colleagues suggested that, based on phylogenic criteria, there are four extant genera of rorquals. They recommend that the genus Balaenoptera be limited to the fin whale, have minke whales fall under the genus Pterobalaena, and have Rorqualus contain the Sei whale, Bryde's whale, Eden's whale (and by extension Rice's whale), the blue whale, and Omura's whale. The gray whale was formerly classified in its own family. The two populations, one in the Sea of Okhotsk and Sea of Japan and the other in eastern Pacific are thought to be genetically and physiologically dissimilar. However, there is some discussion as to whether the gray whale should be classified into its own family, or as a rorqual, with recent studies favoring the latter.
Etymology
The taxonomic name "Mysticeti" () apparently derives from a translation error in early copies of Aristotle's Historia Animalium (in Ancient Greek), in which "" (ho mus to kētos, "the mouse, the whale so called") was mistakenly translated as "" (ho mustikētos, "the Mysticetus"), which D. W. Rice (of the Society for Marine Mammalogy) in assumed was an ironic reference to the animals' great size. An alternate name for the parvorder is "Mystacoceti" (from Greek "mustache" + "whale"), which, although obviously more appropriate and occasionally used in the past, has been superseded by "Mysticeti" (junior synonym).
Mysticetes are also known as baleen whales for their baleen, which they use to sieve plankton and other small organisms from the water. The term "baleen" (Middle English baleyn, ballayne, ballien, bellane, etc.) is an archaic word for "whale", which came from Old French baleine, derived from the Latin word balæna, derived itself from the Ancient Greek φάλλαινα (phállaina).
Right whales got their name because of whalers preferring them over other species; they were essentially the "right whale" to catch.
Differences between families
Rorquals use throat pleats to expand their mouths, which allow them to feed more effectively. However, rorquals need to build up water pressure in order to expand their mouths, leading to a lunge-feeding behavior. Lunge-feeding is where a whale rams a bait ball (a swarm of small fish) at high speed. Rorquals generally have streamlined physiques to reduce drag in the water while doing this.
Balaenids rely on their huge heads, as opposed to the rorquals' throat pleats, to feed effectively. This feeding behavior allows them to grow very big and bulky, without the necessity for a streamlined body. They have callosities, unlike other whales, with the exception of the bowhead whale. Rorquals have a higher proportion of muscle tissue and tend to be negatively buoyant, whereas right whales have a higher proportion of blubber and are positively buoyant. Gray whales are easily distinguished from the other rorquals by their sleet-gray color, dorsal ridges (knuckles on the back), and their gray-white scars left from parasites. As with the other rorquals, their throat pleats increase the capacity of their throats, allowing them to filter larger volumes of water at once. Gray whales are bottom-feeders, meaning they sift through sand to get their food. They usually turn on their sides, scoop up sediment into their mouths and filter out benthic creatures like amphipods, which leave noticeable marks on their heads. The pygmy right whale is easily confused with minke whales because of their similar characteristics, such as their small size, dark gray tops, light gray bottoms, and light eye patches.
Evolutionary history
Molecular phylogeny suggests Mysticeti split from Odontoceti (toothed whales) between 26 and 17 million years ago between the late Oligocene or middle Miocene, but the earliest Mysticeti fossils date to at least 34 million years ago. Their evolutionary link to archaic toothed cetaceans (Archaeoceti) remained unknown until the extinct Janjucetus hunderi was discovered in the early 1990s in Victoria, Australia. While, unlike a modern baleen whale, Janjucetus lacked baleen in its jaw, the anatomy shows sufficient similarity to baleen whales. It appears to have had very limited apparent biosonar capabilities. Its jaw contained teeth, with incisors and canines built for stabbing and molars and premolars built for tearing. These early mysticetes were exceedingly small compared to modern baleen whales, with species like Mammalodon measuring no greater than . It is thought that their size increased with their dependence on baleen. However, the discovery of a skull of the toothed Llanocetus, the second-oldest mysticete, yielded a total length of , indicating filter feeding was not a driving feature in mysticete evolution. The discovery of Janjucetus and others like it suggests that baleen evolution went through several transitional phases. Species like Mammalodon colliveri had little to no baleen, while later species like Aetiocetus weltoni had both baleen and teeth, suggesting they had limited filter feeding capabilities; later genera like Cetotherium had no teeth in their mouth, meaning they were fully dependent on baleen and could only filter feed. However, the 2018 discovery of the toothless Maiabalaena indicates some lineages evolved toothlessness before baleen.
Mystacodon selenensis is the earliest mysticete, dating back to 37 to 33 million years ago (mya) in the Late Eocene, and, like other early toothed mysticetes, or "archaeomysticetes", M. selenensis had heterodont dentition used for suction feeding. Archaeomysticetes from the Oligocene are the Mammalodontidae (Mammalodon and Janjucetus) from Australia. They were small with shortened rostra, and a primitive dental formula (). In baleen whales, it is thought that enlarged mouths adapted for suction feeding evolved before specializations for bulk filter feeding. In the toothed Oligocene mammalodontid Janjucetus, the symphysis is short and the mouth enlarged, the rostrum is wide, and the edges of the maxillae are thin, indicating an adaptation for suction feeding. The aetiocetid Chonecetus still had teeth, but the presence of a groove on the interior side of each mandible indicates the symphysis was elastic, which would have enabled rotation of each mandible, an initial adaptation for bulk feeding like in modern mysticetes.
The first toothless ancestors of baleen whales appeared before the first radiation in the late Oligocene. Eomysticetus and others like it showed no evidence in the skull of echolocation abilities, suggesting they mainly relied on their eyesight for navigation. The eomysticetes had long, flat rostra that lacked teeth and had blowholes located halfway up the dorsal side of the snout. Though the palate is not well-preserved in these specimens, they are thought to have had baleen and been filter feeders. Miocene baleen whales were preyed upon by larger predators like killer sperm whales and megalodon.
The lineages of rorquals and right whales split almost 20 mya. It is unknown where this occurred, but it is generally believed that they, like their descendants, followed plankton migrations. These primitive baleen whales had lost their dentition in favor of baleen, and are believed to have lived on a specialized benthic, plankton, or copepod diet like modern baleen whales. Baleen whales experienced their first radiation in the mid-Miocene. It is thought this radiation was caused by global climate change and major tectonic activity when Antarctica and Australia separated from each other, creating the Antarctic Circumpolar Current. Balaenopterids grew bigger during this time, with species like Balaenoptera sibbaldina perhaps rivaling the blue whale in terms of size, though other studies disagree that any baleen whale grew that large in the Miocene.
The increase in size is likely due to climate change which caused seasonally shifting accumulations of plankton in various parts of the world, necessitating travel over long distances, as well as the ability to feed on large baitballs to make such trips worthwhile. A 2017 analysis of body size based on data from the fossil record and modern baleen whales indicates that the evolution of gigantism in baleen whales occurred rather recently, within the last 3 million years. Before 4.5 million years ago, few baleen whales exceeded in length; the two largest Miocene species were less than in length. The initial evolution of baleen and filter feeding long preceded the evolution of gigantic body size, indicating the evolution of novel feeding mechanisms did not cause the evolution of gigantism. The formation of the Antarctic circumpolar current and its effects on global climate patterns is excluded as being causal for the same reason. Gigantism also was preceded by divergence of different mysticete lineages, meaning multiple lineages arrived at large size independently. It is possible the Plio-Pleistocene increase in seasonally intense upwellings, causing high-prey-density zones, led to gigantism.
Anatomy
Motion
When swimming, baleen whales rely on their flippers for locomotion in a wing-like manner similar to penguins and sea turtles. Flipper movement is continuous. While doing this, baleen whales use their tail fluke to propel themselves forward through vertical motion while using their flippers for steering, much like an otter. Some species leap out of the water, which may allow them to travel faster. Because of their great size, right whales are not flexible or agile like dolphins, and none can move their neck because of the fused cervical vertebrae; this sacrifices speed for stability in the water. The hind legs are enclosed inside the body, and are thought to be vestigial organs. However, a 2014 study suggests that the pelvic bone serves as support for whale genitalia.
Rorquals, needing to build speed to feed, have several adaptions for reducing drag, including a streamlined body; a small dorsal fin, relative to its size; and lack of external ears or long hair. The fin whale is the fastest among baleen whales, having been recorded travelling as fast as , and sustaining a speed of for an extended period. While feeding, the rorqual jaw expands to a volume that can be bigger than the whale itself; to do this, the mouth inflates. The inflation of the mouth causes the cavum ventrale, the throat pleats on the underside stretching to the navel, to expand, increasing the amount of water that the mouth can store. The mandible is connected to the skull by dense fibers and cartilage (fibrocartilage), allowing the jaw to swing open at almost a 90° angle. The mandibular symphysis is also fibrocartilaginous, allowing the jaw to bend which lets in more water. To prevent stretching the mouth too far, rorquals have a sensory organ located in the middle of the jaw to regulate these functions.
External anatomy
Baleen whales have two flippers on the front, near the head. Like all mammals, baleen whales breathe air and must surface periodically to do so. Their nostrils, or blowholes, are situated at the top of the cranium. Baleen whales have two blowholes, as opposed to toothed whales which have one. These paired blowholes are longitudinal slits that converge anteriorly and widen posteriorly, which causes a V-shaped blow. They are surrounded by a fleshy ridge that keeps water away while the whale breathes. The septum that separates the blowholes has two plugs attached to it, making the blowholes water-tight while the whale dives.
Like other mammals, the skin of baleen whales has an epidermis, a dermis, a hypodermis, and connective tissue. The epidermis, the pigmented layer, is thick, along with connective tissue. The epidermis itself is only thick. The dermis, the layer underneath the epidermis, is also thin. The hypodermis, containing blubber, is the thickest part of the skin and functions as a means to conserve heat. Right whales have the thickest hypodermis of any cetacean, averaging , though, as in all whales, it is thinner around openings (such as the blowhole) and limbs. Blubber may also be used to store energy during times of fasting. The connective tissue between the hypodermis and muscles allows only limited movement to occur between them. Unlike toothed whales, baleen whales have small hairs on the top of their head, stretching from the tip of the rostrum to the blowhole, and, in right whales, on the chin. Like other marine mammals, they lack sebaceous and sweat glands.
The baleen of baleen whales are keratinous plates. They are made of a calcified, hard α-keratin material, a fiber-reinforced structure made of intermediate filaments (proteins). The degree of calcification varies between species, with the sei whale having 14.5% hydroxyapatite, a mineral that coats teeth and bones, whereas minke whales have 1–4% hydroxyapatite. In most mammals, keratin structures, such as wool, air-dry, but aquatic whales rely on calcium salts to form on the plates to stiffen them. Baleen plates are attached to the upper jaw and are absent in the mid-jaw, forming two separate combs of baleen. The plates decrease in size as they go further back into the jaw; the largest ones are called the "main baleen plates" and the smallest ones are called the "accessory plates". Accessory plates taper off into small hairs.
Unlike other whales (and most other mammals), the females are larger than the males. Sexual dimorphism is usually reversed, with the males being larger, but the females of all baleen whales are usually five percent larger than males. Sexual dimorphism is also displayed through whale song, notably in humpback whales where the males of the species sing elaborate songs. Male right whales have bigger callosities than female right whales. The males are generally more scarred than females which is thought to be because of aggression during mating season.
Internal systems
The unique lungs of baleen whales are built to collapse under the pressure instead of resisting the pressure which would damage the lungs, enabling some, like the fin whale, to dive to a depth of . The whale lungs are very efficient at extracting oxygen from the air, usually 80%, whereas humans only extract 20% of oxygen from inhaled air. Lung volume is relatively low compared to terrestrial mammals because of the inability of the respiratory tract to hold gas while diving. Doing so may cause serious complications such as embolism. Unlike other mammals, the lungs of baleen whales lack lobes and are more sacculated. Like in humans, the left lung is smaller than the right to make room for the heart. To conserve oxygen, blood is rerouted from pressure-tolerant-tissue to internal organs, and they have a high concentration of myoglobin which allows them to hold their breath longer.
The heart of baleen whales functions similarly to other mammals, with the major difference being the size. The heart can reach , but is still proportional to the whale's size. The muscular wall of the ventricle, which is responsible for pumping blood out of the heart, can be thick. The aorta, an artery, can be thick. Their resting heart rate is 60 to 140 beats per minute (bpm), as opposed to the 60 to 100 bpm in humans. When diving, their heart rate will drop to 4 to 15 bpm to conserve oxygen. Like toothed whales, they have a dense network of blood vessels (rete mirabile) which prevents heat-loss. Like in most mammals, heat is lost in their extremities, so, in baleen whales, warm blood in the arteries is surrounded by veins to prevent heat loss during transport. As well as this, heat inevitably given off by the arteries warms blood in the surrounding veins as it travels back into the core. This is otherwise known as countercurrent exchange. To counteract overheating while in warmer waters, baleen whales reroute blood to the skin to accelerate heat-loss. They have the largest blood corpuscles (red and white blood cells) of any mammal, measuring in diameter, as opposed to human's blood corpuscles.
When sieved from the water, food is swallowed and travels through the esophagus where it enters a three-chambered-stomach. The first compartment is known as the fore-stomach; this is where food gets ground up into an acidic liquid, which is then squirted into the main stomach. Like in humans, the food is mixed with hydrochloric acid and protein-digesting enzymes. Then, the partly digested food is moved into the third stomach, where it meets fat-digesting enzymes, and is then mixed with an alkaline liquid to neutralize the acid from the fore-stomach to prevent damage to the intestinal tract. Their intestinal tract is highly adapted to absorb the most nutrients from food; the walls are folded and contain copious blood vessels, allowing for a greater surface area over which digested food and water can be absorbed. Baleen whales get the water they need from their food; however, the salt content of most of their prey (invertebrates) is similar to that of seawater, whereas the salt content of a whale's blood is considerably lower (three times lower) than that of seawater. The whale kidney is adapted to excreting excess salt; however, while producing urine more concentrated than seawater, it wastes a lot of water which must be replaced.
Baleen whales have a relatively small brain compared to their body mass. Like other mammals, their brain has a large, folded cerebrum, the part of the brain responsible for memory and processing sensory information. Their cerebrum only makes up about 68% of their brain's weight, as opposed to human's 83%. The cerebellum, the part of the brain responsible for balance and coordination, makes up 18% of their brain's weight, compared to 10% in humans, which is probably due to the great degree of control necessary for constantly swimming. Necropsies on the brains of gray whales revealed iron oxide particles, which may allow them to find magnetic north like a compass.
Unlike most animals, whales are conscious breathers. All mammals sleep, but whales cannot afford to become unconscious for long because they may drown. They are believed to exhibit unihemispheric slow-wave sleep, in which they sleep with half of the brain while the other half remains active. This behavior was only documented in toothed whales until footage of a humpback whale sleeping (vertically) was shot in 2014.
It is largely unknown how baleen whales produce sound because of the lack of a melon and vocal cords. Research has found that the larynx had U-shaped folds which are thought to be similar to vocal cords. They are positioned parallel to air flow, as opposed to the perpendicular vocal cords of terrestrial mammals. These may control air flow and cause vibrations. The walls of the larynx are able to contract which may generate sound with support from the arytenoid cartilages. The muscles surrounding the larynx may expel air rapidly or maintain a constant volume while diving.
Senses
The eyes of baleen whales are relatively small for their size and are positioned near the end of the mouth. This is probably because they feed on slow or immobile prey, combined with the fact that most sunlight does not pass , and hence they do not need acute vision. A whale's eye is adapted for seeing both in the euphotic and aphotic zones by increasing or decreasing the pupil's size to prevent damage to the eye. As opposed to land mammals which have a flattened lens, whales have a spherical lens. The retina is surrounded by a reflective layer of cells (tapetum lucidum), which bounces light back at the retina, enhancing eyesight in dark areas. However, light is bent more near the surface of the eye when in air as opposed to water; consequently, they can see much better in the air than in the water. The eyeballs are protected by a thick outer layer to prevent abrasions and an oily fluid (instead of tears) on the surface of the eye. Baleen whales appear to have limited color vision, as they lack S-cones.
The mysticete ear is adapted for hearing underwater, where it can hear sound frequencies as low as 7 Hz and as high as 22 kHz, distinct from odontocetes whose hearing is optimized for ultrasonic frequencies. It is largely unknown how sound is received by baleen whales. Unlike in toothed whales, sound does not pass through the lower jaw. The auditory meatus is blocked by connective tissue and an ear plug, which connects to the eardrum. The inner-ear bones are contained in the tympanic bulla, a bony capsule. However, this is attached to the skull, suggesting that vibrations passing through the bone is important. Sinuses may reflect vibrations towards the cochlea. It is known that when the fluid inside the cochlea is disturbed by vibrations, it triggers sensory hairs which send electric current to the brain, where vibrations are processed into sound.
Baleen whales have a small, yet functional, vomeronasal organ. This allows baleen whales to detect chemicals and pheromones released by their prey. It is thought that 'tasting' the water is important for finding prey and tracking down other whales. They are believed to have an impaired sense of smell due to the lack of the olfactory bulb, but they do have an olfactory tract. Baleen whales have few if any taste buds, suggesting they have lost their sense of taste. They do retain salt-receptor taste-buds suggesting that they can taste saltiness.
Behavior
Migration
Most species of baleen whale migrate long distances from high latitude waters during spring and summer months to more tropical waters during winter months. This migration cycle is repeated annually. The gray whale has the longest recorded migration of any mammal, with one traveling from the Sea of Okhotsk to the Baja Peninsula.
It is thought that plankton blooms dictate where whales migrate. Many baleen whales feed on the massive plankton blooms that occur in the cold, nutrient-rich waters of polar regions during the sunny spring and summer months. Baleen whales generally then migrate to calving grounds in tropical waters during the winter months when plankton populations are low. Migration is hypothesized to benefit calves in a number of ways. Newborns, born with underdeveloped blubber, would likely otherwise be killed by the cold polar temperatures. Migration to warmer waters may also reduce the risk of calves being predated on by killer whales.
Migratory movements may also reflect seasonally shifting patterns of productivity. California blue whales are hypothesized to migrate between dense patches of prey, moving from central California in the summer and fall, to the Gulf of California in the winter, to the central Baja California Pacific coast in spring.
Foraging
All modern mysticetes are obligate filter feeders, using their baleen to strain small prey items (including small fish, krill, copepods, and zooplankton) from seawater. Despite their carnivorous diet, a 2015 study revealed they house gut flora similar to that of terrestrial herbivores. Different kinds of prey are found in different abundances depending on location, and each type of whale is adapted to a specialized way of foraging.
There are two types of feeding behaviors: skim-feeding and lunge-feeding, but some species do both depending on the type and amount of food. Lunge-feeders feed primarily on euphausiids (krill), though some lunge feeders also prey on schools of fish. Skim-feeders, like bowhead whales, feed upon primarily smaller plankton such as copepods. They feed alone or in small groups. Baleen whales get the water they need from their food, and their kidneys excrete excess salt.
The lunge-feeders are the rorquals. To feed, lunge-feeders expand the volume of their jaw to a volume bigger than the original volume of the whale itself. To do this, the mouth inflates, which causes the throat pleats to expand, increasing the amount of water that the mouth can store. Just before they ram the baitball, the jaw swings open at almost a 90° angle and bends which lets in more water. To prevent stretching the mouth too far, rorquals have a sensory organ located in the middle of the jaw to regulate these functions. Then they must decelerate. This process takes a lot of mechanical work and is only energy-effective when used against a large baitball. Lunge feeding is more energy-intensive than skim-feeding due to the acceleration and deceleration required.
The skim-feeders are right whales, gray whales, pygmy right whales, and sei whales (which also lunge feed). To feed, skim-feeders swim with an open mouth, filling it with water and prey. Prey must occur in sufficient numbers to trigger the whale's interest, be within a certain size range so that the baleen plates can filter it, and be slow enough so that it cannot escape. The "skimming" may take place on the surface, underwater, or even at the ocean's bottom, indicated by mud occasionally observed on right whales' bodies. Gray whales feed primarily on the ocean's bottom, feeding on benthic creatures.
Foraging efficiency for both lunge feeding and continuous ram filter feeding is highly dependent upon prey density. The efficiency of a blue whale lunge is approximately 30 times higher at krill densities of than at low krill densities of . Baleen whale have been observed seeking out highly specific areas within the local environment in order to forage at the highest density prey aggregations.
Predation and parasitism
Baleen whales, primarily juveniles and calves, are preyed on by killer whales. It is thought that annual whale migration occurs to protect the calves from the killer whales. There have also been reports of a pod of killer whales attacking and killing an adult bowhead whale, by holding down its flippers, covering the blowhole, and ramming and biting until death. Generally, a mother and calf pair, when faced with the threat of a killer whale pod, will either fight or flee. Fleeing only occurs in species that can swim away quickly, the rorquals. Slower whales must fight the pod alone or with a small family group. There has been one report of a shark attacking and killing a whale calf. This occurred in 2014 during the sardine run when a shiver of dusky sharks attacked a humpback whale calf. Usually, the only shark that will attack a whale is the cookiecutter shark, which leaves a small, non-fatal bite mark.
Many parasites and epibiotics latch onto whales, notably whale lice and whale barnacles. Almost all species of whale lice are specialized towards a certain species of whale, and there can be more than one species per whale. Whale lice eat dead skin, resulting in minor wounds in the skin. Whale louse infestations are especially evident in right whales, where colonies propagate on their callosities. Though not a parasite, whale barnacles latch onto the skin of a whale during their larval stage. However, in doing so it does not harm nor benefit the whale, so their relationship is often labeled as an example of commensalism. Some baleen whales will deliberately rub themselves on substrate to dislodge parasites. Some species of barnacle, such as Conchoderma auritum and whale barnacles, attach to the baleen plates, though this seldom occurs. A species of copepod, Balaenophilus unisetus, inhabits baleen plates of whales. A species of Antarctic diatom, Cocconeis ceticola, forms a film on the skin, which takes a month to develop; this film causes minor damage to the skin. They are also plagued by internal parasites such as stomach worms, cestodes, nematodes, liver flukes, and acanthocephalans.
Reproduction and development
Before reaching adulthood, baleen whales grow at an extraordinary rate. In the blue whale, the largest species, the fetus grows by some per day just before delivery, and by per day during suckling. Before weaning, the calf increases its body weight by and grows from at birth to long. When it reaches sexual maturity after 5–10 years, it will be long and possibly live as long as 80–90 years. Calves are born precocial, needing to be able to swim to the surface at the moment of their birth.
Most rorquals mate in warm waters in winter to give birth almost a year later. A 7-to-11 month lactation period is normally followed by a year of rest before mating starts again. Adults normally start reproducing when 5–10 years old and reach their full length after 20–30 years. In the smallest rorqual, the minke whale, calves are born after a 10-month pregnancy and weaning lasts until it has reached about after 6–7 months. Unusual for a baleen whale, female minkes (and humpbacks) can become pregnant immediately after giving birth; in most species, there is a two-to-three-year calving period. In right whales, the calving interval is usually three years. They grow very rapidly during their first year, after which they hardly increase in size for several years. They reach sexual maturity when long. Baleen whales are K-strategists, meaning they raise one calf at a time, have a long life-expectancy, and a low infant mortality rate. Some 19th century harpoons found in harvested bowheads indicate this species can live more than 100 years. Baleen whales are promiscuous, with none showing pair bonds. They are polygynous, in that a male may mate with more than one female. The scars on male whales suggest they fight for the right to mate with females during breeding season, somewhat similar to lek mating.
Baleen whales have fibroelastic (connective tissue) penises, similar to those of artiodactyls. The tip of the penis, which tapers toward the end, is called the pars intrapraeputialis or terminal cone. The blue whale has the largest penis of any organism on the planet, typically measuring . Accurate measurements of the blue whale are difficult to take because the whale's erect length can only be observed during mating. The penis on a right whale can be up to – the testes, at up to in length, in diameter, and weighing up to , are also the largest of any animal on Earth.
Whale song
All baleen whales use sound for communication and are known to "sing", especially during the breeding season. Blue whales produce the loudest sustained sounds of any animals: their low-frequency (infrasonic, under 20 Hz) moans can last for half a minute, reach almost 190 decibels, and be heard hundreds of kilometers away. Adult male humpbacks produce the longest and most complex songs; sequences of moans, groans, roars, sighs, and chirps sometimes lasting more than ten minutes are repeated for hours. Typically, all humpback males in a population sing the same song over a breeding season, but the songs change slightly between seasons, and males in one population have been observed adapting the song from males of a neighboring population over a few breeding seasons.
Intelligence
Unlike their toothed whale counterparts, baleen whales are hard to study because of their immense size. Intelligence tests such as the mirror test cannot be done because their bulk and lack of body language make a reaction impossible to be definitive. However, studies on the brains of humpback whales revealed spindle cells, which, in humans, control theory of mind. Because of this, it is thought that baleen whales, or at least humpback whales, have consciousness.
Relationship with humans
History of whaling
Whaling by humans has existed since the Stone Age. Ancient whalers used harpoons to spear the bigger animals from boats out at sea. People from Norway started hunting whales around 4,000 years ago, and people from Japan began hunting whales in the Pacific at least as early as that. Whales are typically hunted for their meat and blubber by aboriginal groups; they used baleen for baskets or roofing, and made tools and masks out of bones. The Inuit hunt whales in the Arctic Ocean. The Basques started whaling as early as the 11th century, sailing as far as Newfoundland in the 16th century in search of right whales. 18th and 19th century whalers hunted down whales mainly for their oil, which was used as lamp fuel and a lubricant, and baleen (or whalebone), which was used for items such as corsets and skirt hoops. The most successful whaling nations at this time were the Netherlands, Japan, and the United States.
Commercial whaling was historically important as an industry well throughout the 19th and 20th centuries. Whaling was at that time a sizable European industry with ships from Britain, France, Spain, Denmark, the Netherlands, and Germany, sometimes collaborating to hunt whales in the Arctic. By the early 1790s, whalers, namely the British (Australian) and Americans, started to focus efforts in the South Pacific; in the mid-1900s, over 50,000 humpback whale were taken from the South Pacific. At its height in the 1880s, U.S. profits turned to USD10,000,000, equivalent to US$225,000,000 in 2000. Commonly exploited species included arctic whales such as the gray whale, right whale, and bowhead whale because they were close to the main whaling ports, like New Bedford. After those stocks were depleted, rorquals in the South Pacific were targeted by nearly all whaling organizations; however, they often out-swam whaling vessels. Whaling rorquals was not effective until the harpoon cannon was invented in the late 1860s. Whaling basically stopped when stocks of all species were depleted to a point that they could not be harvested on a commercial scale. Whaling was controlled in 1982 when the International Whaling Commission (IWC) placed a moratorium setting catch limits to protect species from dying out from over-exploitation, and eventually banned it:
Conservation and management issues
As of 2021, the International Union for Conservation of Nature (IUCN) recognizes 15 mysticete species (while not yet officially recognizing Rice's whale as a species, it still gives it a conservation status as a distinct population segment). Two species—the North Atlantic right whale (with only around 366 individuals left) and Rice's whale (with less than 100 individuals left)—are considered critically endangered. Three more are classified as endangered (the North Pacific right whale, the blue whale, and the sei whale), one as vulnerable (the fin whale), one as near-threatened (Antarctic minke whale), and one as data deficient (Omura's whale). Species that live in polar habitats are vulnerable to the effects of ongoing climate change, particularly declines in sea ice, as well as ocean acidification.
The whale-watching industry and anti-whaling advocates argue that whaling catches "friendly" whales that are curious about boats, as these whales are the easiest to catch. This analysis claims that once the economic benefits of hotels, restaurants and other tourist amenities are considered, hunting whales is a net economic loss. This argument is particularly contentious in Iceland, as it has among the most-developed whale-watching operations in the world and the hunting of minke whales resumed in August 2003. Brazil, Argentina and South Africa argue that whale watching is a growing billion-dollar industry that provides more revenue than commercial whaling would provide. Peru, Uruguay, Australia, and New Zealand also support proposals to permanently forbid whaling south of the Equator, as Solor (an island of Indonesia) is the only place of the Southern Hemisphere that takes whales. Anti-whaling groups, such as the International Fund for Animal Welfare (IFAW), claim that countries which support a pro-whaling stance are damaging their economies by driving away anti-whaling tourists.
Commercial whaling was historically important for the world economy. All species were exploited, and as one type's stock depleted, another type was targeted. The scale of whale harvesting decreased substantially through the 1960s as all whale stocks had been depleted, and practically stopped in 1988 after the International Whaling Commission placed a moratorium which banned whaling for commercial use. Several species that were commercially exploited have rebounded in numbers; for example, gray whales may be as numerous as they were prior to whaling, making it the first marine mammal to be taken off the endangered species list. The Southern right whale was hunted to near extinction in the mid-to-late 20th century, with only a small (unknown) population around Antarctica. Because of international protection, the Southern right whale's population has been growing 7% annually since 1970. Conversely, the eastern stock of North Atlantic right whale was extirpated from much of its former range, which stretched from the coast of North Africa to the North Sea and Iceland; it is thought that the entire stock consists of only ten individuals, making the eastern stock functionally extinct.
Baleen whales continue to be harvested. Only three nations take whales: Iceland, Norway, and Japan. All these nations are part of the IWC, with Norway and Iceland rejecting the moratorium and continuing commercial whaling. Japan, being part of the IWC, whales under the Scientific Permit stated in Article VIII in the Convention for the Regulation of Whaling, which allows the taking of whales for scientific research. Japan has had two main research programs: the Joint Aquatic Resources Permit Application (JARPA) and the Japanese Research Program in the North (JARPN). JARPN is focused in the North Pacific and JARPA around the Antarctic. JARPA mainly caught Antarctic minke whales, catching nearly 7,000; to a far lesser extent, they also caught fin whales. Animal-rights activist groups, such as the Greenpeace, object to Japan's scientific whaling, with some calling it a substitute for commercial whaling. In 2014, the International Court of Justice (the UN judicial branch) banned the taking of whales for any purpose in the Southern Ocean Whale Sanctuary; however, Japan refuses to stop whaling and has only promised to cut their annual catches by a third (around 300 whales per year).
Baleen whales can also be affected by humans in more indirect ways. For species like the North Atlantic right whale, which migrates through some of the world's busiest shipping lanes, the biggest threat is from being struck by ships. The Lloyd's mirror effect results in low frequency propeller sounds not being discernible near the surface, where most accidents occur. Combined with spreading and acoustic shadowing effects, the result is that the whale is unable to hear an approaching vessel before it has been run over or entrapped by the hydrodynamic forces of the vessel's passage. A 2014 study noted that a lower vessel speed correlated with lower collision rates. The ever-increasing amount of ocean noise, including sonar, drowns out the vocalizations produced by whales, notably in the blue whale which produces the loudest vocalization, which makes it harder for them to communicate. Blue whales stop producing foraging D calls once a mid-frequency sonar is activated, even though the sonar frequency range (1–8 kHz) far exceeds their sound production range (25–100 Hz).
Poisoning from toxic substances such as polychlorinated biphenyl (PCB) is generally low because of their low trophic level. However, oil spills can be a significant threat, especially to small populations; the already endangered Rice's whale was likely devastated by the Deepwater Horizon oil spill, with some estimates indicating a decline of up to 22% in the species.
Some baleen whales can become victims of bycatch, which is especially serious for North Atlantic right whales considering their small number. Right whales feed with a wide-open mouth, risking entanglement in any rope or net fixed in the water column. The rope wraps around their upper jaw, flippers and tail. Some are able to escape, but others remain entangled. If observers notice, they can be successfully disentangled, but others die over a period of months. Other whales, such as humpback whales, can also be entangled.
In captivity
Baleen whales have rarely been kept in captivity. Their large size and appetite make them expensive creatures to maintain. Pools of proper size would also be very expensive to build. For example, a single gray whale calf would need to eat of fish per day, and the pool would have to accommodate the calf, along with ample room to swim. Only gray whales have survived being kept in captivity for over a year. The first gray whale, which was captured in Scammon's Lagoon, Baja California Sur, in 1965, was named Gigi and died two months later from an infection. The second gray whale, which was captured in 1971 from the same lagoon, was named Gigi II and was released a year later after becoming too big. The last gray whale, J.J., beached itself in Marina del Rey, California, where it was rushed to SeaWorld San Diego and, after 14 months, was released because it got too big to take care of. Reaching and , J.J. was the largest creature to be kept in captivity.
The Mito Aquarium in Numazu, Shizuoka, Japan, housed three minke whales in the nearby bay enclosed by nets. One survived for three months, another (a calf) survived for two weeks, and another was kept for over a month before breaking through the nets.
| Biology and health sciences | Baleen whales | Animals |
231064 | https://en.wikipedia.org/wiki/Terrestrial%20television | Terrestrial television | Terrestrial television, or over-the-air television (OTA) is a type of television broadcasting in which the content is transmitted via radio waves from the terrestrial (Earth-based) transmitter of a TV station to a TV receiver having an antenna. The term terrestrial is more common in Europe and Latin America, while in Canada and the United States it is called over-the-air or simply broadcast. This type of TV broadcast is distinguished from newer technologies, such as satellite television (direct broadcast satellite or DBS television), in which the signal is transmitted to the receiver from an overhead satellite; cable television, in which the signal is carried to the receiver through a cable; and Internet Protocol television, in which the signal is received over an Internet stream or on a network utilizing the Internet Protocol. Terrestrial television stations broadcast on television channels with frequencies between about 52 and 600 MHz in the VHF and UHF bands. Since radio waves in these bands travel by line of sight, reception is generally limited by the visual horizon to distances of , although under better conditions and with tropospheric ducting, signals can sometimes be received hundreds of kilometers distant.
Terrestrial television was the first technology used for television broadcasting. The BBC began broadcasting in 1929 and by 1930 many radio stations had a regular schedule of experimental television programmes. However, these early experimental systems had insufficient picture quality to attract the public, due to their mechanical scan technology, and television did not become widespread until after World War II with the advent of electronic scan television technology. The television broadcasting business followed the model of radio networks, with local television stations in cities and towns affiliated with television networks, either commercial (in the US) or government-controlled (in Europe), which provided content. Television broadcasts were in greyscale (called black and white) until the transition to color television in the 1960s.
There was no other method of television delivery until the 1950s with the beginnings of cable television and community antenna television (CATV). CATV was, initially, only a re-broadcast of over-the-air signals. With the widespread adoption of cable across the United States in the 1970s and 1980s, viewing of terrestrial television broadcasts has been in decline; in 2018, it was estimated that about 14% of US households used an antenna. However, in certain other regions terrestrial television continue to be the preferred method of receiving television, and it is estimated by Deloitte as of 2020 that at least 1.6 billion people in the world receive at least some television using these means. The largest market is thought to be Indonesia, where 250 million people watch through terrestrial.
By 2019, over-the-top media service (OTT) which is streamed via the internet had become a common alternative.
Analog
Europe
Following the ST61 conference, UHF frequencies were first used in the UK in 1964 with the introduction of BBC2. In the UK, VHF channels were kept on the old 405-line system, while UHF was used solely for 625-line broadcasts (which later used PAL color). Television broadcasting in the 405-line system continued after the introduction of four analog programs in the UHF bands until the last 405-line transmitters were switched off on January 6, 1985. VHF Band III was used in other countries around Europe for PAL broadcasts until the planned phase-out and switch over to digital television.
The success of analog terrestrial television across Europe varied from country to country. Although each country had rights to a certain number of frequencies by virtue of the ST61 plan, not all of them were brought into service.
Americas
The first National Television System Committee standard was introduced in 1941. This standard defined a transmission scheme for a black-and-white picture with 525 lines of vertical resolution at 60 fields per second. In the early 1950s, this standard was superseded by a backward-compatible standard for color television. The NTSC standard was exclusively being used in the Americas as well as Japan until the introduction of digital terrestrial television (DTT). While Mexico has ended all its analog television broadcasts and the United States and Canada have shut down nearly all of their analog TV stations, the NTSC standard continues to be used in the rest of Latin American countries except for Argentina, Paraguay and Uruguay where PAL-N standard is used while testing their DTT platform.
In the late 1990s and early 2000s, the Advanced Television Systems Committee developed the ATSC standard for digital high-definition terrestrial transmission. This standard was eventually adopted by many American countries, including the United States, Canada, Dominican Republic, Mexico, Argentina, El Salvador, Guatemala and Honduras; however, the four latter countries reversed their decision in favor of ISDB-Tb.
The Pan-American terrestrial television operates on analog channels 2 through 6 (VHF-low band, 54 to 88 MHz, known as band I in Europe), 7 through 13 (VHF-high band, 174 to 216 MHz, known as band III elsewhere), and 14 through 51 (UHF television band, 470 to 698 MHz, elsewhere bands IV and V). Unlike with analog transmission, ATSC channel numbers do not correspond to radio frequencies. Instead, a virtual channel is defined as part of the ATSC stream metadata so that a station can transmit on any frequency but still show the same channel number. Additionally, free-to-air television repeaters and signal boosters can be used to rebroadcast a terrestrial television signal using an otherwise unused channel to cover areas with marginal reception. (see Pan-American television frequencies for frequency allocation charts)
Analog television channels 2 through 6, 7 through 13, and 14 through 51 are only used for LPTV translator stations in the United States. Channels 52 through 69 are still used by some existing stations, but these channels must be vacated if telecommunications companies notify the stations to vacate that signal spectrum. By convention, broadcast television signals are transmitted with horizontal polarization.
Asia
Terrestrial television broadcast in Asia started as early as 1939 in Japan through a series of experiments done by NHK Broadcasting Institute of Technology. However, these experiments were interrupted by the beginning of the World War II in the Pacific. On February 1, 1953, NHK (Japan Broadcasting Corporation) began broadcasting. On August 28, 1953, Nippon TV (Nippon Television Network Corporation), the first commercial television broadcaster in Asia was launched. Meanwhile, in the Philippines, Alto Broadcasting System (now ABS-CBN Corporation), the first commercial television broadcaster in Southeast Asia, launched its first commercial terrestrial television station DZAQ-TV on October 23, 1953, with the help of Radio Corporation of America (RCA).
Digital
By the mid-1990s, the interest in digital television across Europe was such the CEPT convened the "Chester '97" conference to agree on means by which digital television could be inserted into the ST61 frequency plan.
The introduction of digital terrestrial television in the late 1990s and early years of the 21st century led the ITU to call a Regional Radiocommunication Conference to abrogate the ST61 plan and to put a new plan for DTT broadcasting only in its place.
In December 2005, the European Union decided to cease all analog audio and analog video television transmissions by 2012 and switch all terrestrial television broadcasting to digital audio and digital video (all EU countries have agreed on using DVB-T). The Netherlands completed the transition in December 2006, and some EU member states decided to complete their switchover as early as 2008 (Sweden), and (Denmark) in 2009. While the UK began to switch off analog broadcasts, region by region, in late 2007, it was not completed until 24 October 2012. Norway ceased all analog television transmissions on 1 December 2009. Two member states (not specified in the announcement) expressed concerns that they might not be able to proceed to the switchover by 2012 due to technical limitations; the rest of the EU member states had stopped analog television transmissions by the end 2012.
Many countries are developing and evaluating digital terrestrial television systems.
Australia has adopted the DVB-T standards and the government's industry regulator, the Australian Communications and Media Authority, has mandated that all analog transmissions will cease by 2012. Mandated digital conversion started early in 2009 with a graduated program. The first centre to experience analog switch-off was the remote Victorian regional town of Mildura, in 2010. The government supplied underprivileged houses across the nation with free digital set-top converter boxes in order to minimize conversion disruption. Australia's major free-to-air television networks were all granted digital transmission licenses and are each required to broadcast at least one high-definition and one standard-definition channel into all of their markets.
In North America, a specification laid out by the ATSC has become the standard for digital terrestrial television. In the United States, the Federal Communications Commission (FCC) set the final deadline for the switch-off of analog service for 12 June 2009. All television receivers must now include a DTT tuner using ATSC. In Canada, the Canadian Radio-television and Telecommunications Commission (CRTC) set 31 August 2011 as the date that terrestrial analog transmission service ceased in metropolitan areas and provincial capitals.
In Mexico, the Federal Telecommunications Institute (IFT) discontinued the use of analog terrestrial television on 31 December 2015.
| Technology | Broadcasting | null |
231137 | https://en.wikipedia.org/wiki/Earthquake%20prediction | Earthquake prediction | Earthquake prediction is a branch of the science of seismology concerned with the specification of the time, location, and magnitude of future earthquakes within stated limits, and particularly "the determination of parameters for the next strong earthquake to occur in a region". Earthquake prediction is sometimes distinguished from earthquake forecasting, which can be defined as the probabilistic assessment of general earthquake hazard, including the frequency and magnitude of damaging earthquakes in a given area over years or decades.
Prediction can be further distinguished from earthquake warning systems, which, upon detection of an earthquake, provide a real-time warning of seconds to neighboring regions that might be affected.
In the 1970s, scientists were optimistic that a practical method for predicting earthquakes would soon be found, but by the 1990s continuing failure led many to question whether it was even possible. Demonstrably successful predictions of large earthquakes have not occurred, and the few claims of success are controversial. For example, the most famous claim of a successful prediction is that alleged for the 1975 Haicheng earthquake. A later study said that there was no valid short-term prediction. Extensive searches have reported many possible earthquake precursors, but, so far, such precursors have not been reliably identified across significant spatial and temporal scales. While part of the scientific community hold that, taking into account non-seismic precursors and given enough resources to study them extensively, prediction might be possible, most scientists are pessimistic and some maintain that earthquake prediction is inherently impossible.
Evaluating earthquake predictions
Predictions are deemed significant if they can be shown to be successful beyond random chance. Therefore, methods of statistical hypothesis testing are used to determine the probability that an earthquake such as is predicted would happen anyway (the null hypothesis). The predictions are then evaluated by testing whether they correlate with actual earthquakes better than the null hypothesis.
In many instances, however, the statistical nature of earthquake occurrence is not simply homogeneous. Clustering occurs in both space and time. In southern California about 6% of M≥3.0 earthquakes are "followed by an earthquake of larger magnitude within 5 days and 10 km." In central Italy 9.5% of M≥3.0 earthquakes are followed by a larger event within 48 hours and 30 km. While such statistics are not satisfactory for purposes of prediction (giving ten to twenty false alarms for each successful prediction) they will skew the results of any analysis that assumes that earthquakes occur randomly in time, for example, as realized from a Poisson process. It has been shown that a "naive" method based solely on clustering can successfully predict about 5% of earthquakes; "far better than 'chance'".
As the purpose of short-term prediction is to enable emergency measures to reduce death and destruction, failure to give warning of a major earthquake, that does occur, or at least an adequate evaluation of the hazard, can result in legal liability, or even political purging. For example, it has been reported that members of the Chinese Academy of Sciences were purged for "having ignored scientific predictions of the disastrous Tangshan earthquake of summer 1976." Following the 2009 L'Aquila Earthquake, seven scientists and technicians in Italy were convicted of manslaughter, but not so much for failing to predict the earthquake, where some 300 people died, as for giving undue assurance to the populace – one victim called it "anaesthetizing" – that there would not be a serious earthquake, and therefore no need to take precautions. But warning of an earthquake that does not occur also incurs a cost: not only the cost of the emergency measures themselves, but of civil and economic disruption. False alarms, including alarms that are canceled, also undermine the credibility, and thereby the effectiveness, of future warnings. In 1999 it was reported that China was introducing "tough regulations intended to stamp out 'false' earthquake warnings, in order to prevent panic and mass evacuation of cities triggered by forecasts of major tremors." This was prompted by "more than 30 unofficial earthquake warnings ... in the past three years, none of which has been accurate." The acceptable trade-off between missed quakes and false alarms depends on the societal valuation of these outcomes. The rate of occurrence of both must be considered when evaluating any prediction method.
In a 1997 study of the cost-benefit ratio of earthquake prediction research in Greece, Stathis Stiros suggested that even a (hypothetical) excellent prediction method would be of questionable social utility, because "organized evacuation of urban centers is unlikely to be successfully accomplished", while "panic and other undesirable side-effects can also be anticipated." He found that earthquakes kill less than ten people per year in Greece (on average), and that most of those fatalities occurred in large buildings with identifiable structural issues. Therefore, Stiros stated that it would be much more cost-effective to focus efforts on identifying and upgrading unsafe buildings. Since the death toll on Greek highways is more than 2300 per year on average, he argued that more lives would also be saved if Greece's entire budget for earthquake prediction had been used for street and highway safety instead.
Prediction methods
Earthquake prediction is an immature scienceit has not yet led to a successful prediction of an earthquake from first physical principles. Research into methods of prediction therefore focus on empirical analysis, with two general approaches: either identifying distinctive precursors to earthquakes, or identifying some kind of geophysical trend or pattern in seismicity that might precede a large earthquake. Precursor methods are pursued largely because of their potential utility for short-term earthquake prediction or forecasting, while 'trend' methods are generally thought to be useful for forecasting, long term prediction (10 to 100 years time scale) or intermediate term prediction (1 to 10 years time scale).
Precursors
An earthquake precursor is an anomalous phenomenon that might give effective warning of an impending earthquake. Reports of these – though generally recognized as such only after the event – number in the thousands, some dating back to antiquity. There have been around 400 reports of possible precursors in scientific literature, of roughly twenty different types, running the gamut from aeronomy to zoology. None have been found to be reliable for the purposes of earthquake prediction.
In the early 1990, the IASPEI solicited nominations for a Preliminary List of Significant Precursors. Forty nominations were made, of which five were selected as possible significant precursors, with two of those based on a single observation each.
After a critical review of the scientific literature, the International Commission on Earthquake Forecasting for Civil Protection (ICEF) concluded in 2011 there was "considerable room for methodological improvements in this type of research." In particular, many cases of reported precursors are contradictory, lack a measure of amplitude, or are generally unsuitable for a rigorous statistical evaluation. Published results are biased towards positive results, and so the rate of false negatives (earthquake but no precursory signal) is unclear.
Animal behavior
After an earthquake has already begun, pressure waves (P waves) travel twice as fast as the more damaging shear waves (s waves). Typically not noticed by humans, some animals may notice the smaller vibrations that arrive a few to a few dozen seconds before the main shaking, and become alarmed or exhibit other unusual behavior. Seismometers can also detect P waves, and the timing difference is exploited by electronic earthquake warning systems to provide humans with a few seconds to move to a safer location.
A review of scientific studies available as of 2018 covering over 130 species found insufficient evidence to show that animals could provide warning of earthquakes hours, days, or weeks in advance. Statistical correlations suggest some reported unusual animal behavior is due to smaller earthquakes (foreshocks) that sometimes precede a large quake, which if small enough may go unnoticed by people. Foreshocks may also cause groundwater changes or release gases that can be detected by animals. Foreshocks are also detected by seismometers, and have long been studied as potential predictors, but without success (see #Seismicity patterns). Seismologists have not found evidence of medium-term physical or chemical changes that predict earthquakes which animals might be sensing.
Anecdotal reports of strange animal behavior before earthquakes have been recorded for thousands of years. Some unusual animal behavior may be mistakenly attributed to a near-future earthquake. The flashbulb memory effect causes unremarkable details to become more memorable and more significant when associated with an emotionally powerful event such as an earthquake. Even the vast majority of scientific reports in the 2018 review did not include observations showing that animals did not act unusually when there was not an earthquake about to happen, meaning the behavior was not established to be predictive.
Most researchers investigating animal prediction of earthquakes are in China and Japan. Most scientific observations have come from the 2010 Canterbury earthquake in New Zealand, the 1984 Nagano earthquake in Japan, and the 2009 L'Aquila earthquake in Italy.
Animals known to be magnetoreceptive might be able to detect electromagnetic waves in the ultra low frequency and extremely low frequency ranges that reach the surface of the Earth before an earthquake, causing odd behavior. These electromagnetic waves could also cause air ionization, water oxidation and possible water toxification which other animals could detect.
Dilatancy–diffusion
In the 1970s the dilatancy–diffusion hypothesis was highly regarded as providing a physical basis for various phenomena seen as possible earthquake precursors. It was based on "solid and repeatable evidence" from laboratory experiments that highly stressed crystalline rock experienced a change in volume, or dilatancy, which causes changes in other characteristics, such as seismic velocity and electrical resistivity, and even large-scale uplifts of topography. It was believed this happened in a 'preparatory phase' just prior to the earthquake, and that suitable monitoring could therefore warn of an impending quake.
Detection of variations in the relative velocities of the primary and secondary seismic waves – expressed as Vp/Vs – as they passed through a certain zone was the basis for predicting the 1973 Blue Mountain Lake (NY) and 1974 Riverside (CA) quake. Although these predictions were informal and even trivial, their apparent success was seen as confirmation of both dilatancy and the existence of a preparatory process, leading to what were subsequently called "wildly over-optimistic statements" that successful earthquake prediction "appears to be on the verge of practical reality."
However, many studies questioned these results, and the hypothesis eventually languished. Subsequent study showed it "failed for several reasons, largely associated with the validity of the assumptions on which it was based", including the assumption that laboratory results can be scaled up to the real world. Another factor was the bias of retrospective selection of criteria. Other studies have shown dilatancy to be so negligible that concluded: "The concept of a large-scale 'preparation zone' indicating the likely magnitude of a future event, remains as ethereal as the ether that went undetected in the Michelson–Morley experiment."
Changes in Vp/Vs
Vp is the symbol for the velocity of a seismic "P" (primary or pressure) wave passing through rock, while Vs is the symbol for the velocity of the "S" (secondary or shear) wave. Small-scale laboratory experiments have shown that the ratio of these two velocities – represented as Vp/Vs – changes when rock is near the point of fracturing. In the 1970s it was considered a likely breakthrough when Russian seismologists reported observing such changes (later discounted.) in the region of a subsequent earthquake. This effect, as well as other possible precursors, has been attributed to dilatancy, where rock stressed to near its breaking point expands (dilates) slightly.
Study of this phenomenon near Blue Mountain Lake in New York State led to a successful albeit informal prediction in 1973, and it was credited for predicting the 1974 Riverside (CA) quake. However, additional successes have not followed, and it has been suggested that these predictions were a fluke. A Vp/Vs anomaly was the basis of a 1976 prediction of a M 5.5 to 6.5 earthquake near Los Angeles, which failed to occur. Other studies relying on quarry blasts (more precise, and repeatable) found no such variations, while an analysis of two earthquakes in California found that the variations reported were more likely caused by other factors, including retrospective selection of data. noted that reports of significant velocity changes have ceased since about 1980.
Radon emissions
Most rock contains small amounts of gases that can be isotopically distinguished from the normal atmospheric gases. There are reports of spikes in the concentrations of such gases prior to a major earthquake; this has been attributed to release due to pre-seismic stress or fracturing of the rock. One of these gases is radon, produced by radioactive decay of the trace amounts of uranium present in most rock. Radon is potentially useful as an earthquake predictor because it is radioactive and thus easily detected, and its short half-life (3.8 days) makes radon levels sensitive to short-term fluctuations.
A 2009 compilation listed 125 reports of changes in radon emissions prior to 86 earthquakes since 1966. The International Commission on Earthquake Forecasting for Civil Protection (ICEF) however found in its 2011 critical review that the earthquakes with which these changes are supposedly linked were up to a thousand kilometers away, months later, and at all magnitudes. In some cases the anomalies were observed at a distant site, but not at closer sites. The ICEF found "no significant correlation".
Electromagnetic anomalies
Observations of electromagnetic disturbances and their attribution to the earthquake failure process go back as far as the Great Lisbon earthquake of 1755, but practically all such observations prior to the mid-1960s are invalid because the instruments used were sensitive to physical movement. Since then various anomalous electrical, electric-resistive, and magnetic phenomena have been attributed to precursory stress and strain changes that precede earthquakes, raising hopes for finding a reliable earthquake precursor. While a handful of researchers have gained much attention with either theories of how such phenomena might be generated, claims of having observed such phenomena prior to an earthquake, no such phenomena has been shown to be an actual precursor.
A 2011 review by the International Commission on Earthquake Forecasting for Civil Protection (ICEF) found the "most convincing" electromagnetic precursors to be ultra low frequency magnetic anomalies, such as the Corralitos event (discussed below) recorded before the 1989 Loma Prieta earthquake. However, it is now believed that observation was a system malfunction. Study of the closely monitored 2004 Parkfield earthquake found no evidence of precursory electromagnetic signals of any type; further study showed that earthquakes with magnitudes less than 5 do not produce significant transient signals. The ICEF considered the search for useful precursors to have been unsuccessful.
VAN seismic electric signals
The most touted, and most criticized, claim of an electromagnetic precursor is the VAN method of physics professors Panayiotis Varotsos, Kessar Alexopoulos and Konstantine Nomicos (VAN) of the University of Athens. In a 1981 paper they claimed that by measuring geoelectric voltages – what they called "seismic electric signals" (SES) – they could predict earthquakes.
In 1984, they claimed there was a "one-to-one correspondence" between SES and earthquakes – that is, that "every sizable EQ is preceded by an SES and inversely every SES is always followed by an EQ the magnitude and the epicenter of which can be reliably predicted" – the SES appearing between 6 and 115 hours before the earthquake. As proof of their method they claimed a series of successful predictions.
Although their report was "saluted by some as a major breakthrough", among seismologists it was greeted by a "wave of generalized skepticism". In 1996, a paper VAN submitted to the journal Geophysical Research Letters was given an unprecedented public peer-review by a broad group of reviewers, with the paper and reviews published in a special issue; the majority of reviewers found the methods of VAN to be flawed. Additional criticism was raised the same year in a public debate between some of the principals.
A primary criticism was that the method is geophysically implausible and scientifically unsound. Additional objections included the demonstrable falsity of the claimed one-to-one relationship of earthquakes and SES, the unlikelihood of a precursory process generating signals stronger than any observed from the actual earthquakes, and the very strong likelihood that the signals were man-made. Further work in Greece has tracked SES-like "anomalous transient electric signals" back to specific human sources, and found that such signals are not excluded by the criteria used by VAN to identify SES. More recent work, by employing modern methods of statistical physics, i.e., detrended fluctuation analysis (DFA), multifractal DFA and wavelet transform revealed that SES are clearly distinguished from signals produced by man made sources.
The validity of the VAN method, and therefore the predictive significance of SES, was based primarily on the empirical claim of demonstrated predictive success. Numerous weaknesses have been uncovered in the VAN methodology, and in 2011 the International Commission on Earthquake Forecasting for Civil Protection concluded that the prediction capability claimed by VAN could not be validated. Most seismologists consider VAN to have been "resoundingly debunked". On the other hand, the Section "Earthquake Precursors and Prediction" of "Encyclopedia of Solid Earth Geophysics: part of "Encyclopedia of Earth Sciences Series" (Springer 2011) ends as follows (just before its summary): "it has recently been shown that by analyzing time-series in a newly introduced time domain "natural time", the approach to the critical state can be clearly identified [Sarlis et al. 2008]. This way, they appear to have succeeded in shortening the lead-time of VAN prediction to only a few days [Uyeda and Kamogawa 2008]. This means, seismic data may play an amazing role in short term precursor when combined with SES data".
Since 2001, the VAN group has introduced a concept they call "natural time", applied to the analysis of their precursors. Initially it is applied on SES to distinguish them from noise and relate them to a possible impending earthquake. In case of verification (classification as "SES activity"), natural time analysis is additionally applied to the general subsequent seismicity of the area associated with the SES activity, in order to improve the time parameter of the prediction. The method treats earthquake onset as a critical phenomenon. A review of the updated VAN method in 2020 says that it suffers from an abundance of false positives and is therefore not usable as a prediction protocol. VAN group answered by pinpointing misunderstandings in the specific reasoning.
Corralitos anomaly
Probably the most celebrated seismo-electromagnetic event ever, and one of the most frequently cited examples of a possible earthquake precursor, is the 1989 Corralitos anomaly. In the month prior to the 1989 Loma Prieta earthquake, measurements of the Earth's magnetic field at ultra-low frequencies by a magnetometer in Corralitos, California, just 7 km from the epicenter of the impending earthquake, started showing anomalous increases in amplitude. Just three hours before the quake, the measurements soared to about thirty times greater than normal, with amplitudes tapering off after the quake. Such amplitudes had not been seen in two years of operation, nor in a similar instrument located 54 km away. To many people such apparent locality in time and space suggested an association with the earthquake.
Additional magnetometers were subsequently deployed across northern and southern California, but after ten years and several large earthquakes, similar signals have not been observed. More recent studies have cast doubt on the connection, attributing the Corralitos signals to either unrelated magnetic disturbance or, even more simply, to sensor-system malfunction.
Freund physics
In his investigations of crystalline physics, Friedemann Freund found that water molecules embedded in rock can dissociate into ions if the rock is under intense stress. The resulting charge carriers can generate battery currents under certain conditions. Freund suggested that perhaps these currents could be responsible for earthquake precursors such as electromagnetic radiation, earthquake lights and disturbances of the plasma in the ionosphere. The study of such currents and interactions is known as "Freund physics".
Most seismologists reject Freund's suggestion that stress-generated signals can be detected and put to use as precursors, for a number of reasons. First, it is believed that stress does not accumulate rapidly before a major earthquake, and thus there is no reason to expect large currents to be rapidly generated. Secondly, seismologists have extensively searched for statistically reliable electrical precursors, using sophisticated instrumentation, and have not identified any such precursors. And thirdly, water in the Earth's crust would cause any generated currents to be absorbed before reaching the surface.
Disturbance of the daily cycle of the ionosphere
The ionosphere usually develops its lower D layer during the day, while at night this layer disappears as the plasma there turns to gas. During the night, the F layer of the ionosphere remains formed, in higher altitude than D layer. A waveguide for low HF radio frequencies up to 10 MHz is formed during the night (skywave propagation) as the F layer reflects these waves back to the Earth. The skywave is lost during the day, as the D layer absorbs these waves.
Tectonic stresses in the Earth's crust are claimed to cause waves of electric charges that travel to the surface of the Earth and affect the ionosphere. ULF* recordings of the daily cycle of the ionosphere indicate that the usual cycle could be disturbed a few days before a shallow strong earthquake. When the disturbance occurs, it is observed that either the D layer is lost during the day resulting to ionosphere elevation and skywave formation or the D layer appears at night resulting to lower of the ionosphere and hence absence of skywave.
Science centers have developed a network of VLF transmitters and receivers on a global scale that detect changes in skywave. Each receiver is also daisy transmitter for distances of 1000–10,000 kilometers and is operating at different frequencies within the network. The general area under excitation can be determined depending on the density of the network. It was shown on the other hand that global extreme events like magnetic storms or solar flares and local extreme events in the same VLF path like another earthquake or a volcano eruption that occur in near time with the earthquake under evaluation make it difficult or impossible to relate changes in skywave to the earthquake of interest.
In 2017, an article in the Journal of Geophysical Research showed that the relationship between ionospheric anomalies and large seismic events (M≥6.0) occurring globally from 2000 to 2014 was based on the presence of solar weather. When the solar data are removed from the time series, the correlation is no longer statistically significant. A subsequent article in Physics of the Earth and Planetary Interiors in 2020 shows that solar weather and ionospheric disturbances are a potential cause to trigger large earthquakes based on this statistical relationship. The proposed mechanism is electromagnetic induction from the ionosphere to the fault zone. Fault fluids are conductive, and can produce telluric currents at depth. The resulting change in the local magnetic field in the fault triggers dissolution of minerals and weakens the rock, while also potentially changing the groundwater chemistry and level. After the seismic event, different minerals may be precipitated thus changing groundwater chemistry and level again. This process of mineral dissolution and precipitation before and after an earthquake has been observed in Iceland. This model makes sense of the ionospheric, seismic and groundwater data.
Satellite observation of the expected ground temperature declination
One way of detecting the mobility of tectonic stresses is to detect locally elevated temperatures on the surface of the crust measured by satellites. During the evaluation process, the background of daily variation and noise due to atmospheric disturbances and human activities are removed before visualizing the concentration of trends in the wider area of a fault. This method has been experimentally applied since 1995.
In a newer approach to explain the phenomenon, NASA's Friedmann Freund has proposed that the infrared radiation captured by the satellites is not due to a real increase in the surface temperature of the crust. According to this version the emission is a result of the quantum excitation that occurs at the chemical re-bonding of positive charge carriers (holes) which are traveling from the deepest layers to the surface of the crust at a speed of 200 meters per second. The electric charge arises as a result of increasing tectonic stresses as the time of the earthquake approaches. This emission extends superficially up to 500 x 500 square kilometers for very large events and stops almost immediately after the earthquake.
Trends
Instead of watching for anomalous phenomena that might be precursory signs of an impending earthquake, other approaches to predicting earthquakes look for trends or patterns that lead to an earthquake. As these trends may be complex and involve many variables, advanced statistical techniques are often needed to understand them, therefore these are sometimes called statistical methods. These approaches also tend to be more probabilistic, and to have larger time periods, and so merge into earthquake forecasting.
Nowcasting
Earthquake nowcasting, suggested in 2016 is the estimate of the current dynamic state of a seismological system, based on natural time introduced in 2001. It differs from forecasting which aims to estimate the probability of a future event but it is also considered a potential base for forecasting. Nowcasting calculations produce the "earthquake potential score", an estimation of the current level of seismic progress. Typical applications are: great global earthquakes and tsunamis, aftershocks and induced seismicity, induced seismicity at gas fields, seismic risk to global megacities, studying of clustering of large global earthquakes, etc.
Elastic rebound
Even the stiffest of rock is not perfectly rigid. Given a large force (such as between two immense tectonic plates moving past each other) the Earth's crust will bend or deform. According to the elastic rebound theory of , eventually the deformation (strain) becomes great enough that something breaks, usually at an existing fault. Slippage along the break (an earthquake) allows the rock on each side to rebound to a less deformed state. In the process energy is released in various forms, including seismic waves. The cycle of tectonic force being accumulated in elastic deformation and released in a sudden rebound is then repeated. As the displacement from a single earthquake ranges from less than a meter to around 10 meters (for an M 8 quake), the demonstrated existence of large strike-slip displacements of hundreds of miles shows the existence of a long running earthquake cycle.
Characteristic earthquakes
The most studied earthquake faults (such as the Nankai megathrust, the Wasatch Fault, and the San Andreas Fault) appear to have distinct segments. The characteristic earthquake model postulates that earthquakes are generally constrained within these segments. As the lengths and other properties of the segments are fixed, earthquakes that rupture the entire fault should have similar characteristics. These include the maximum magnitude (which is limited by the length of the rupture), and the amount of accumulated strain needed to rupture the fault segment. Since continuous plate motions cause the strain to accumulate steadily, seismic activity on a given segment should be dominated by earthquakes of similar characteristics that recur at somewhat regular intervals. For a given fault segment, identifying these characteristic earthquakes and timing their recurrence rate (or conversely return period) should therefore inform us about the next rupture; this is the approach generally used in forecasting seismic hazard. UCERF3 is a notable example of such a forecast, prepared for the state of California. Return periods are also used for forecasting other rare events, such as cyclones and floods, and assume that future frequency will be similar to observed frequency to date.
The idea of characteristic earthquakes was the basis of the Parkfield prediction: fairly similar earthquakes in 1857, 1881, 1901, 1922, 1934, and 1966 suggested a pattern of breaks every 21.9 years, with a standard deviation of ±3.1 years. Extrapolation from the 1966 event led to a prediction of an earthquake around 1988, or before 1993 at the latest (at the 95% confidence interval). The appeal of such a method is that the prediction is derived entirely from the trend, which supposedly accounts for the unknown and possibly unknowable earthquake physics and fault parameters. However, in the Parkfield case the predicted earthquake did not occur until 2004, a decade late. This seriously undercuts the claim that earthquakes at Parkfield are quasi-periodic, and suggests the individual events differ sufficiently in other respects to question whether they have distinct characteristics in common.
The failure of the Parkfield prediction has raised doubt as to the validity of the characteristic earthquake model itself. Some studies have questioned the various assumptions, including the key one that earthquakes are constrained within segments, and suggested that the "characteristic earthquakes" may be an artifact of selection bias and the shortness of seismological records (relative to earthquake cycles). Other studies have considered whether other factors need to be considered, such as the age of the fault. Whether earthquake ruptures are more generally constrained within a segment (as is often seen), or break past segment boundaries (also seen), has a direct bearing on the degree of earthquake hazard: earthquakes are larger where multiple segments break, but in relieving more strain they will happen less often.
Seismic gaps
At the contact where two tectonic plates slip past each other every section must eventually slip, as (in the long-term) none get left behind. But they do not all slip at the same time; different sections will be at different stages in the cycle of strain (deformation) accumulation and sudden rebound. In the seismic gap model the "next big quake" should be expected not in the segments where recent seismicity has relieved the strain, but in the intervening gaps where the unrelieved strain is the greatest. This model has an intuitive appeal; it is used in long-term forecasting, and was the basis of a series of circum-Pacific (Pacific Rim) forecasts in 1979 and 1989–1991.
However, some underlying assumptions about seismic gaps are now known to be incorrect. A close examination suggests that "there may be no information in seismic gaps about the time of occurrence or the magnitude of the next large event in the region"; statistical tests of the circum-Pacific forecasts shows that the seismic gap model "did not forecast large earthquakes well". Another study concluded that a long quiet period did not increase earthquake potential.
Seismicity patterns
Various heuristically derived algorithms have been developed for predicting earthquakes. Probably the most widely known is the M8 family of algorithms (including the RTP method) developed under the leadership of Vladimir Keilis-Borok. M8 issues a "Time of Increased Probability" (TIP) alarm for a large earthquake of a specified magnitude upon observing certain patterns of smaller earthquakes. TIPs generally cover large areas (up to a thousand kilometers across) for up to five years. Such large parameters have made M8 controversial, as it is hard to determine whether any hits that happened were skillfully predicted, or only the result of chance.
M8 gained considerable attention when the 2003 San Simeon and Hokkaido earthquakes occurred within a TIP. In 1999, Keilis-Borok's group published a claim to have achieved statistically significant intermediate-term results using their M8 and MSc models, as far as world-wide large earthquakes are regarded. However, Geller et al. are skeptical of prediction claims over any period shorter than 30 years. A widely publicized TIP for an M 6.4 quake in Southern California in 2004 was not fulfilled, nor two other lesser known TIPs. A deep study of the RTP method in 2008 found that out of some twenty alarms only two could be considered hits (and one of those had a 60% chance of happening anyway). It concluded that "RTP is not significantly different from a naïve method of guessing based on the historical rates [of] seismicity."
Accelerating moment release (AMR, "moment" being a measurement of seismic energy), also known as time-to-failure analysis, or accelerating seismic moment release (ASMR), is based on observations that foreshock activity prior to a major earthquake not only increased, but increased at an exponential rate. In other words, a plot of the cumulative number of foreshocks gets steeper just before the main shock.
Following formulation by into a testable hypothesis, and a number of positive reports, AMR seemed promising despite several problems. Known issues included not being detected for all locations and events, and the difficulty of projecting an accurate occurrence time when the tail end of the curve gets steep. But rigorous testing has shown that apparent AMR trends likely result from how data fitting is done, and failing to account for spatiotemporal clustering of earthquakes. The AMR trends are therefore statistically insignificant. Interest in AMR (as judged by the number of peer-reviewed papers) has fallen off since 2004.
Machine learning
Rouet-Leduc et al. (2019) reported having successfully trained a regression random forest on acoustic time series data capable of identifying a signal emitted from fault zones that forecasts fault failure. Rouet-Leduc et al. (2019) suggested that the identified signal, previously assumed to be statistical noise, reflects the increasing emission of energy before its sudden release during a slip event. Rouet-Leduc et al. (2019) further postulated that their approach could bound fault failure times and lead to the identification of other unknown signals. Due to the rarity of the most catastrophic earthquakes, acquiring representative data remains problematic. In response, Rouet-Leduc et al. (2019) have conjectured that their model would not need to train on data from catastrophic earthquakes, since further research has shown the seismic patterns of interest to be similar in smaller earthquakes.
Deep learning has also been applied to earthquake prediction. Although Bath's law and Omori's law describe the magnitude of earthquake aftershocks and their time-varying properties, the prediction of the "spatial distribution of aftershocks" remains an open research problem. Using the Theano and TensorFlow software libraries, DeVries et al. (2018) trained a neural network that achieved higher accuracy in the prediction of spatial distributions of earthquake aftershocks than the previously established methodology of Coulomb failure stress change. Notably, DeVries et al. (2018) reported that their model made no "assumptions about receiver plane orientation or geometry" and heavily weighted the change in shear stress, "sum of the absolute values of the independent components of the stress-change tensor," and the von Mises yield criterion. DeVries et al. (2018) postulated that the reliance of their model on these physical quantities indicated that they might "control earthquake triggering during the most active part of the seismic cycle." For validation testing, DeVries et al. (2018) reserved 10% of positive training earthquake data samples and an equal quantity of randomly chosen negative samples.
Arnaud Mignan and Marco Broccardo have similarly analyzed the application of artificial neural networks to earthquake prediction. They found in a review of literature that earthquake prediction research utilizing artificial neural networks has gravitated towards more sophisticated models amidst increased interest in the area. They also found that neural networks utilized in earthquake prediction with notable success rates were matched in performance by simpler models. They further addressed the issues of acquiring appropriate data for training neural networks to predict earthquakes, writing that the "structured, tabulated nature of earthquake catalogues" makes transparent machine learning models more desirable than artificial neural networks.
EMP induced seismicity
High energy electromagnetic pulses can induce earthquakes within 2–6 days after the emission by EMP generators. It has been proposed that strong EM impacts could control seismicity, as the seismicity dynamics that follow appear to be a lot more regular than usual.
Notable predictions
These are predictions, or claims of predictions, that are notable either scientifically or because of public notoriety, and claim a scientific or quasi-scientific basis. As many predictions are held confidentially, or published in obscure locations, and become notable only when they are claimed, there may be a selection bias in that hits get more attention than misses. The predictions listed here are discussed in Hough's book and Geller's paper.
1975: Haicheng, China
The M 7.3 1975 Haicheng earthquake is the most widely cited "success" of earthquake prediction. The ostensible story is that study of seismic activity in the region led the Chinese authorities to issue a medium-term prediction in June 1974, and the political authorities therefore ordered various measures taken, including enforced evacuation of homes, construction of "simple outdoor structures", and showing of movies out-of-doors. The quake, striking at 19:36, was powerful enough to destroy or badly damage about half of the homes. However, the "effective preventative measures taken" were said to have kept the death toll under 300 in an area with population of about 1.6 million, where otherwise tens of thousands of fatalities might have been expected.
However, although a major earthquake occurred, there has been some skepticism about the narrative of measures taken on the basis of a timely prediction. This event occurred during the Cultural Revolution, when "belief in earthquake prediction was made an element of ideological orthodoxy that distinguished the true party liners from right wing deviationists". Recordkeeping was disordered, making it difficult to verify details, including whether there was any ordered evacuation. The method used for either the medium-term or short-term predictions (other than "Chairman Mao's revolutionary line") has not been specified. The evacuation may have been spontaneous, following the strong (M 4.7) foreshock that occurred the day before.
A 2006 study that had access to an extensive range of records found that the predictions were flawed. "In particular, there was no official short-term prediction, although such a prediction was made by individual scientists." Also: "it was the foreshocks alone that triggered the final decisions of warning and evacuation". They estimated that 2,041 lives were lost. That more did not die was attributed to a number of fortuitous circumstances, including earthquake education in the previous months (prompted by elevated seismic activity), local initiative, timing (occurring when people were neither working nor asleep), and local style of construction. The authors conclude that, while unsatisfactory as a prediction, "it was an attempt to predict a major earthquake that for the first time did not end up with practical failure."
1981: Lima, Peru (Brady)
In 1976, Brian Brady, a physicist, then at the U.S. Bureau of Mines, where he had studied how rocks fracture, "concluded a series of four articles on the theory of earthquakes with the deduction that strain building in the subduction zone [off-shore of Peru] might result in an earthquake of large magnitude within a period of seven to fourteen years from mid November 1974." In an internal memo written in June 1978 he narrowed the time window to "October to November, 1981", with a main shock in the range of 9.2±0.2. In a 1980 memo he was reported as specifying "mid-September 1980". This was discussed at a scientific seminar in San Juan, Argentina, in October 1980, where Brady's colleague, W. Spence, presented a paper. Brady and Spence then met with government officials from the U.S. and Peru on 29 October, and "forecast a series of large magnitude earthquakes in the second half of 1981." This prediction became widely known in Peru, following what the U.S. embassy described as "sensational first page headlines carried in most Lima dailies" on January 26, 1981.
On 27 January 1981, after reviewing the Brady-Spence prediction, the U.S. National Earthquake Prediction Evaluation Council (NEPEC) announced it was "unconvinced of the scientific validity" of the prediction, and had been "shown nothing in the observed seismicity data, or in the theory insofar as presented, that lends substance to the predicted times, locations, and magnitudes of the earthquakes." It went on to say that while there was a probability of major earthquakes at the predicted times, that probability was low, and recommend that "the prediction not be given serious consideration."
Unfazed, Brady subsequently revised his forecast, stating there would be at least three earthquakes on or about July 6, August 18 and September 24, 1981, leading one USGS official to complain: "If he is allowed to continue to play this game ... he will eventually get a hit and his theories will be considered valid by many."
On June 28 (the date most widely taken as the date of the first predicted earthquake), it was reported that: "the population of Lima passed a quiet Sunday". The headline on one Peruvian newspaper: "NO PASÓ NADA" ("Nothing happened").
In July Brady formally withdrew his prediction on the grounds that prerequisite seismic activity had not occurred. Economic losses due to reduced tourism during this episode has been roughly estimated at one hundred million dollars.
1985–1993: Parkfield, U.S. (Bakun-Lindh)
The "Parkfield earthquake prediction experiment" was the most heralded scientific earthquake prediction ever. It was based on an observation that the Parkfield segment of the San Andreas Fault breaks regularly with a moderate earthquake of about M 6 every several decades: 1857, 1881, 1901, 1922, 1934, and 1966. More particularly, pointed out that, if the 1934 quake is excluded, these occur every 22 years, ±4.3 years. Counting from 1966, they predicted a 95% chance that the next earthquake would hit around 1988, or 1993 at the latest. The National Earthquake Prediction Evaluation Council (NEPEC) evaluated this, and concurred. The U.S. Geological Survey and the State of California therefore established one of the "most sophisticated and densest nets of monitoring instruments in the world", in part to identify any precursors when the quake came. Confidence was high enough that detailed plans were made for alerting emergency authorities if there were signs an earthquake was imminent. In the words of The Economist: "never has an ambush been more carefully laid for such an event."
Year 1993 came, and passed, without fulfillment. Eventually there was an M 6.0 earthquake on the Parkfield segment of the fault, on 28 September 2004, but without forewarning or obvious precursors. While the experiment in catching an earthquake is considered by many scientists to have been successful, the prediction was unsuccessful in that the eventual event was a decade late.
1983–1995: Greece (VAN)
In 1981, the "VAN" group, headed by Panayiotis Varotsos, said that they found a relationship between earthquakes and 'seismic electric signals' (SES). In 1984 they presented a table of 23 earthquakes from 19 January 1983 to 19 September 1983, of which they claimed to have successfully predicted 18 earthquakes. Other lists followed, such as their 1991 claim of predicting six out of seven earthquakes with ≥ 5.5 in the period of 1 April 1987 through 10 August 1989, or five out of seven earthquakes with ≥ 5.3 in the overlapping period of 15 May 1988 to 10 August 1989, In 1996 they published a "Summary of all Predictions issued from January 1st, 1987 to June 15, 1995", amounting to 94 predictions. Matching this against a list of "All earthquakes with MS(ATH)" and within geographical bounds including most of Greece, they come up with a list of 14 earthquakes they should have predicted. Here they claim ten successes, for a success rate of 70%.
The VAN predictions have been criticized on various grounds, including being geophysically implausible, "vague and ambiguous", failing to satisfy prediction criteria, and retroactive adjustment of parameters. A critical review of 14 cases where VAN claimed 10 successes showed only one case where an earthquake occurred within the prediction parameters. The VAN predictions not only fail to do better than chance, but show "a much better association with the events which occurred before them", according to Mulargia and Gasperini. Other early reviews found that the VAN results, when evaluated by definite parameters, were statistically significant. Both positive and negative views on VAN predictions from this period were summarized in the 1996 book A Critical Review of VAN edited by Sir James Lighthill and in a debate issue presented by the journal Geophysical Research Letters that was focused on the statistical significance of the VAN method. VAN had the opportunity to reply to their critics in those review publications. In 2011, the ICEF reviewed the 1996 debate, and concluded that the optimistic SES prediction capability claimed by VAN could not be validated. In 2013, the SES activities were found to be coincident with the minima of the fluctuations of the order parameter of seismicity, which have been shown to be statistically significant precursors by employing the event coincidence analysis.
A crucial issue is the large and often indeterminate parameters of the predictions, such that some critics say these are not predictions, and should not be recognized as such. Much of the controversy with VAN arises from this failure to adequately specify these parameters. Some of their telegrams include predictions of two distinct earthquake events, such as (typically) one earthquake predicted at 300 km "NW" of Athens, and another at 240 km "W", "with 5,3 and 5,8", with no time limit. The time parameter estimation was introduced in VAN Method by means of natural time in 2001.
VAN has disputed the 'pessimistic' conclusions of their critics, but the critics have not relented. It was suggested that VAN failed to account for clustering of earthquakes, or that they interpreted their data differently during periods of greater seismic activity.
VAN has been criticized on several occasions for causing public panic and widespread unrest. This has been exacerbated by the broadness of their predictions, which cover large areas of Greece (up to 240 kilometers across, and often pairs of areas), much larger than the areas actually affected by earthquakes of the magnitudes predicted (usually several tens of kilometers across). Magnitudes are similarly broad: a predicted magnitude of "6.0" represents a range from a benign magnitude 5.3 to a broadly destructive 6.7. Coupled with indeterminate time windows of a month or more, such predictions "cannot be practically utilized" to determine an appropriate level of preparedness, whether to curtail usual societal functioning, or even to issue public warnings.
2008: Greece (VAN)
After 2006, VAN claim that all alarms related to SES activity have been made public by posting at arxiv.org. Such SES activity is evaluated using a new method they call 'natural time'. One such report was posted on Feb. 1, 2008, two weeks before the strongest earthquake in Greece during the period 1983–2011. This earthquake occurred on February 14, 2008, with magnitude (Mw) 6.9. VAN's report was also described in an article in the newspaper Ethnos on Feb. 10, 2008. However, Gerassimos Papadopoulos commented that the VAN reports were confusing and ambiguous, and that "none of the claims for successful VAN predictions is justified." A reply to this comment, which insisted on the prediction's accuracy, was published in the same issue.
1989: Loma Prieta, U.S.
The 1989 Loma Prieta earthquake (epicenter in the Santa Cruz Mountains northwest of San Juan Bautista, California) caused significant damage in the San Francisco Bay Area of California. The United States Geological Survey (USGS) reportedly claimed, twelve hours after the event, that it had "forecast" this earthquake in a report the previous year. USGS staff subsequently claimed this quake had been "anticipated"; various other claims of prediction have also been made.
Ruth Harris () reviewed 18 papers (with 26 forecasts) dating from 1910 "that variously offer or relate to scientific forecasts of the 1989 Loma Prieta earthquake." (In this case no distinction is made between a forecast, which is limited to a probabilistic estimate of an earthquake happening over some time period, and a more specific prediction.) None of these forecasts can be rigorously tested due to lack of specificity, and where a forecast does bracket the correct time and location, the window was so broad (e.g., covering the greater part of California for five years) as to lose any value as a prediction. Predictions that came close (but given a probability of only 30%) had ten- or twenty-year windows.
One debated prediction came from the M8 algorithm used by Keilis-Borok and associates in four forecasts. The first of these forecasts missed both magnitude (M 7.5) and time (a five-year window from 1 January 1984, to 31 December 1988). They did get the location, by including most of California and half of Nevada. A subsequent revision, presented to the NEPEC, extended the time window to 1 July 1992, and reduced the location to only central California; the magnitude remained the same. A figure they presented had two more revisions, for M ≥ 7.0 quakes in central California. The five-year time window for one ended in July 1989, and so missed the Loma Prieta event; the second revision extended to 1990, and so included Loma Prieta.
When discussing success or failure of prediction for the Loma Prieta earthquake, some scientists argue that it did not occur on the San Andreas Fault (the focus of most of the forecasts), and involved dip-slip (vertical) movement rather than strike-slip (horizontal) movement, and so was not predicted.
Other scientists argue that it did occur in the San Andreas Fault zone, and released much of the strain accumulated since the 1906 San Francisco earthquake; therefore several of the forecasts were correct. Hough states that "most seismologists" do not believe this quake was predicted "per se". In a strict sense there were no predictions, only forecasts, which were only partially successful.
Iben Browning claimed to have predicted the Loma Prieta event, but (as will be seen in the next section) this claim has been rejected.
1990: New Madrid, U.S. (Browning)
Iben Browning (a scientist with a Ph.D. degree in zoology and training as a biophysicist, but no experience in geology, geophysics, or seismology) was an "independent business consultant" who forecast long-term climate trends for businesses. He supported the idea (scientifically unproven) that volcanoes and earthquakes are more likely to be triggered when the tidal force of the Sun and the Moon coincide to exert maximum stress on the Earth's crust (syzygy). Having calculated when these tidal forces maximize, Browning then "projected" what areas were most at risk for a large earthquake. An area he mentioned frequently was the New Madrid seismic zone at the southeast corner of the state of Missouri, the site of three very large earthquakes in 1811–12, which he coupled with the date of 3 December 1990.
Browning's reputation and perceived credibility were boosted when he claimed in various promotional flyers and advertisements to have predicted (among various other events) the Loma Prieta earthquake of 17 October 1989. The National Earthquake Prediction Evaluation Council (NEPEC) formed an Ad Hoc Working Group (AHWG) to evaluate Browning's prediction. Its report (issued 18 October 1990) specifically rejected the claim of a successful prediction of the Loma Prieta earthquake. A transcript of his talk in San Francisco on 10 October showed he had said: "there will probably be several earthquakes around the world, Richter 6+, and there may be a volcano or two" – which, on a global scale, is about average for a week – with no mention of any earthquake in California.
Though the AHWG report disproved both Browning's claims of prior success and the basis of his "projection", it made little impact after a year of continued claims of a successful prediction. Browning's prediction received the support of geophysicist David Stewart, and the tacit endorsement of many public authorities in their preparations for a major disaster, all of which was amplified by massive exposure in the news media. Nothing happened on 3 December, and Browning died of a heart attack seven months later.
2004 and 2005: Southern California, U.S. (Keilis-Borok)
The M8 algorithm (developed under the leadership of Vladimir Keilis-Borok at UCLA) gained respect by the apparently successful predictions of the 2003 San Simeon and Hokkaido earthquakes. Great interest was therefore generated by the prediction in early 2004 of a M ≥ 6.4 earthquake to occur somewhere within an area of southern California of approximately 12,000 sq. miles, on or before 5 September 2004. In evaluating this prediction the California Earthquake Prediction Evaluation Council (CEPEC) noted that this method had not yet made enough predictions for statistical validation, and was sensitive to input assumptions. It therefore concluded that no "special public policy actions" were warranted, though it reminded all Californians "of the significant seismic hazards throughout the state." The predicted earthquake did not occur.
A very similar prediction was made for an earthquake on or before 14 August 2005, in approximately the same area of southern California. The CEPEC's evaluation and recommendation were essentially the same, this time noting that the previous prediction and two others had not been fulfilled. This prediction also failed.
2009: L'Aquila, Italy (Giuliani)
At 03:32 on 6 April 2009, the Abruzzo region of central Italy was rocked by a magnitude M 6.3 earthquake. In the city of L'Aquila and surrounding area around 60,000 buildings collapsed or were seriously damaged, resulting in 308 deaths and 67,500 people left homeless. Around the same time, it was reported that Giampaolo Giuliani had predicted the earthquake, had tried to warn the public, but had been muzzled by the Italian government.
Giampaolo Giuliani was a laboratory technician at the Laboratori Nazionali del Gran Sasso. As a hobby he had for some years been monitoring radon using instruments he had designed and built. Prior to the L'Aquila earthquake he was unknown to the scientific community, and had not published any scientific work. He had been interviewed on 24 March by an Italian-language blog, Donne Democratiche, about a swarm of low-level earthquakes in the Abruzzo region that had started the previous December. He said that this swarm was normal and would diminish by the end of March. On 30 March, L'Aquila was struck by a magnitude 4.0 temblor, the largest to date.
On 27 March Giuliani warned the mayor of L'Aquila there could be an earthquake within 24 hours, and an earthquake M~2.3 occurred. On 29 March he made a second prediction. He telephoned the mayor of the town of Sulmona, about 55 kilometers southeast of L'Aquila, to expect a "damaging" – or even "catastrophic" – earthquake within 6 to 24 hours. Loudspeaker vans were used to warn the inhabitants of Sulmona to evacuate, with consequential panic. No quake ensued and Giuliano was cited for inciting public alarm and enjoined from making future public predictions.
After the L'Aquila event Giuliani claimed that he had found alarming rises in radon levels just hours before. He said he had warned relatives, friends and colleagues on the evening before the earthquake hit. He was subsequently interviewed by the International Commission on Earthquake Forecasting for Civil Protection, which found that Giuliani had not transmitted a valid prediction of the mainshock to the civil authorities before its occurrence.
Difficulty or impossibility
As the preceding examples show, the record of earthquake prediction has been disappointing. The optimism of the 1970s that routine prediction of earthquakes would be "soon", perhaps within ten years, was coming up disappointingly short by the 1990s, and many scientists began wondering why. By 1997 it was being positively stated that earthquakes can not be predicted, which led to a notable debate in 1999 on whether prediction of individual earthquakes is a realistic scientific goal.
Earthquake prediction may have failed only because it is "fiendishly difficult" and still beyond the current competency of science. Despite the confident announcement four decades ago that seismology was "on the verge" of making reliable predictions, there may yet be an underestimation of the difficulties. As early as 1978 it was reported that earthquake rupture might be complicated by "heterogeneous distribution of mechanical properties along the fault", and in 1986 that geometrical irregularities in the fault surface "appear to exert major controls on the starting and stopping of ruptures". Another study attributed significant differences in fault behavior to the maturity of the fault. These kinds of complexities are not reflected in current prediction methods.
Seismology may even yet lack an adequate grasp of its most central concept, elastic rebound theory. A simulation that explored assumptions regarding the distribution of slip found results "not in agreement with the classical view of the elastic rebound theory". (This was attributed to details of fault heterogeneity not accounted for in the theory.)
Earthquake prediction may be intrinsically impossible. In 1997, it has been argued that the Earth is in a state of self-organized criticality "where any small earthquake has some probability of cascading into a large event". It has also been argued on decision-theoretic grounds that "prediction of major earthquakes is, in any practical sense, impossible." In 2021, a multitude of authors from a variety of universities and research institutes studying the China Seismo-Electromagnetic Satellite reported that the claims based on self-organized criticality stating that at any moment any small earthquake can eventually cascade to a large event, do not stand in view of the results obtained to date by natural time analysis.
That earthquake prediction might be intrinsically impossible has been strongly disputed, but the best disproof of impossibility – effective earthquake prediction – has yet to be demonstrated.
| Physical sciences | Seismology | Earth science |
231145 | https://en.wikipedia.org/wiki/Krill | Krill | Krill (Euphausiids) (: krill) are small and exclusively marine crustaceans of the order Euphausiacea, found in all the world's oceans. The name "krill" comes from the Norwegian word , meaning "small fry of fish", which is also often attributed to species of fish.
Krill are considered an important trophic level connection near the bottom of the food chain. They feed on phytoplankton and, to a lesser extent, zooplankton, and are also the main source of food for many larger animals. In the Southern Ocean, one species, the Antarctic krill, makes up an estimated biomass of around 379 million tonnes, making it among the species with the largest total biomass. Over half of this biomass is eaten by whales, seals, penguins, seabirds, squid, and fish each year. Most krill species display large daily vertical migrations, providing food for predators near the surface at night and in deeper waters during the day.
Krill are fished commercially in the Southern Ocean and in the waters around Japan. The total global harvest amounts to 150,000–200,000 tonnes annually, mostly from the Scotia Sea. Most krill catch is used for aquaculture and aquarium feeds, as bait in sport fishing, or in the pharmaceutical industry. Krill are also used for human consumption in several countries. They are known as in Japan and as camarones in Spain and the Philippines. In the Philippines, they are also called alamang and are used to make a salty paste called bagoong.
Krill are also the main prey of baleen whales, including the blue whale.
Taxonomy
Krill belong to the large arthropod subphylum, the Crustacea. The most familiar and largest group of crustaceans, the class Malacostraca, includes the superorder Eucarida comprising the three orders, Euphausiacea (krill), Decapoda (shrimp, prawns, lobsters, crabs), and the planktonic Amphionidacea.
The order Euphausiacea comprises two families. The more abundant Euphausiidae contains 10 different genera with a total of 85 species. Of these, the genus Euphausia is the largest, with 31 species. The lesser-known family, the Bentheuphausiidae, has only one species, Bentheuphausia amblyops, a bathypelagic krill living in deep waters below . It is considered the most primitive extant krill species.
Well-known species of the Euphausiidae of commercial krill fisheries include Antarctic krill (Euphausia superba), Pacific krill (E. pacifica) and Northern krill (Meganyctiphanes norvegica).
Phylogeny
, the order Euphausiacea is believed to be monophyletic due to several unique conserved morphological characteristics (autapomorphy) such as its naked filamentous gills and thin thoracopods and by molecular studies.
There have been many theories of the location of the order Euphausiacea. Since the first description of Thysanopode tricuspide by Henri Milne-Edwards in 1830, the similarity of their biramous thoracopods had led zoologists to group euphausiids and Mysidacea in the order Schizopoda, which was split by Johan Erik Vesti Boas in 1883 into two separate orders. Later, William Thomas Calman (1904) ranked the Mysidacea in the superorder Peracarida and euphausiids in the superorder Eucarida, although even up to the 1930s the order Schizopoda was advocated. It was later also proposed that order Euphausiacea should be grouped with the Penaeidae (family of prawns) in the Decapoda based on developmental similarities, as noted by Robert Gurney and Isabella Gordon. The reason for this debate is that krill share some morphological features of decapods and others of mysids.
Molecular studies have not unambiguously grouped them, possibly due to the paucity of key rare species such as Bentheuphausia amblyops in krill and Amphionides reynaudii in Eucarida. One study supports the monophyly of Eucarida (with basal Mysida), another groups Euphausiacea with Mysida (the Schizopoda), while yet another groups Euphausiacea with Hoplocarida.
Timeline
No extant fossil can be unequivocally assigned to Euphausiacea. Some extinct eumalacostracan taxa have been thought to be euphausiaceans such as Anthracophausia, Crangopsis—now assigned to the Aeschronectida (Hoplocarida)—and Palaeomysis. All dating of speciation events were estimated by molecular clock methods, which placed the last common ancestor of the krill family Euphausiidae (order Euphausiacea minus Bentheuphausia amblyops) to have lived in the Lower Cretaceous about .
Distribution
Krill occur worldwide in all oceans, although many individual species have endemic or neritic (i.e., coastal) distributions. Bentheuphausia amblyops, a bathypelagic species, has a cosmopolitan distribution within its deep-sea habitat.
Species of the genus Thysanoessa occur in both Atlantic and Pacific oceans. The Pacific is home to Euphausia pacifica. Northern krill occur across the Atlantic from the Mediterranean Sea northward.
Species with neritic distributions include the four species of the genus Nyctiphanes. They are highly abundant along the upwelling regions of the California, Humboldt, Benguela, and Canarias current systems. Another species having only neritic distribution is E. crystallorophias, which is endemic to the Antarctic coastline.
Species with endemic distributions include Nyctiphanes capensis, which occurs only in the Benguela current, E. mucronata in the Humboldt current, and the six Euphausia species native to the Southern Ocean.
In the Antarctic, seven species are known, one in genus Thysanoessa (T. macrura) and six in Euphausia. The Antarctic krill (Euphausia superba) commonly lives at depths reaching , whereas ice krill (Euphausia crystallorophias) reach depth of , though they commonly inhabit depths of at most . Krill perform Diel Vertical Migrations (DVM) in large swarms, and acoustic data has shown these migrations to go up to 400 metres in depth. Both are found at latitudes south of 55° S, with E. crystallorophias dominating south of 74° S and in regions of pack ice. Other species known in the Southern Ocean are E. frigida, E. longirostris, E. triacantha and E. vallentini.
Anatomy and morphology
Krill are crustaceans and, like all crustaceans, they have a chitinous exoskeleton. They have anatomy similar to a standard decapod with their bodies made up of three parts: the cephalothorax is composed of the head and the thorax, which are fused, and the abdomen, which bears the ten swimming appendages, and the tail fan. This outer shell of krill is transparent in most species.
Krill feature intricate compound eyes. Some species adapt to different lighting conditions through the use of screening pigments.
They have two antennae and several pairs of thoracic legs called pereiopods or thoracopods, so named because they are attached to the thorax. Their number varies among genera and species. These thoracic legs include feeding legs and grooming legs.
Krill are probably the sister clade of decapods because all species have five pairs of swimming legs called "swimmerets" in common with the latter, very similar to those of a lobster or freshwater crayfish.
In spite of having ten swimmerets, otherwise known as pleopods, krill cannot be considered decapods. They lack any true ground-based legs due to all their pereiopods having been converted into grooming and auxiliary feeding legs. In Decapoda, there are ten functioning pereiopods, giving them their name; whereas here there are no remaining locomotive pereiopods. Nor are there consistently ten pereiopods at all.
Most krill are about long as adults. A few species grow to sizes on the order of . The largest krill species, Thysanopoda cornuta, lives deep in the open ocean. Krill can be easily distinguished from other crustaceans such as true shrimp by their externally visible gills.
Except for Bentheuphausia amblyops, krill are bioluminescent animals having organs called photophores that can emit light. The light is generated by an enzyme-catalysed chemiluminescence reaction, wherein a luciferin (a kind of pigment) is activated by a luciferase enzyme. Studies indicate that the luciferin of many krill species is a fluorescent tetrapyrrole similar but not identical to dinoflagellate luciferin and that the krill probably do not produce this substance themselves but acquire it as part of their diet, which contains dinoflagellates. Krill photophores are complex organs with lenses and focusing abilities, and can be rotated by muscles. The precise function of these organs is as yet unknown; possibilities include mating, social interaction or orientation and as a form of counter-illumination camouflage to compensate their shadow against overhead ambient light.
Ecology
Feeding
Many krill are filter feeders: their frontmost appendages, the thoracopods, form very fine combs with which they can filter out their food from the water. These filters can be very fine in species (such as Euphausia spp.) that feed primarily on phytoplankton, in particular on diatoms, which are unicellular algae. Krill are mostly omnivorous, although a few species are carnivorous, preying on small zooplankton and fish larvae.
Krill are an important element of the aquatic food chain. Krill convert the primary production of their prey into a form suitable for consumption by larger animals that cannot feed directly on the minuscule algae. Northern krill and some other species have a relatively small filtering basket and actively hunt copepods and larger zooplankton.
Predation
Many animals feed on krill, ranging from smaller animals like fish or penguins to larger ones like seals and baleen whales.
Disturbances of an ecosystem resulting in a decline in the krill population can have far-reaching effects. During a coccolithophore bloom in the Bering Sea in 1998, for instance, the diatom concentration dropped in the affected area. Krill cannot feed on the smaller coccolithophores, and consequently the krill population (mainly E. pacifica) in that region declined sharply. This in turn affected other species: the shearwater population dropped. The incident was thought to have been one reason salmon did not spawn that season.
Several single-celled endoparasitoidic ciliates of the genus Collinia can infect species of krill and devastate affected populations. Such diseases were reported for Thysanoessa inermis in the Bering Sea and also for E. pacifica, Thysanoessa spinifera, and T. gregaria off the North American Pacific coast. Some ectoparasites of the family Dajidae (epicaridean isopods) afflict krill (and also shrimp and mysids); one such parasite is Oculophryxus bicaulis, which was found on the krill Stylocheiron affine and S. longicorne. It attaches itself to the animal's eyestalk and sucks blood from its head; it apparently inhibits the host's reproduction, as none of the afflicted animals reached maturity.
Climate change poses another threat to krill populations.
Plastics
Preliminary research indicates krill can digest microplastics under in diameter, breaking them down and excreting them back into the environment in smaller form.
Life history and behavior
The life cycle of krill is relatively well understood, despite minor variations in detail from species to species. After krill hatch, they experience several larval stages—nauplius, pseudometanauplius, metanauplius, calyptopsis, and furcilia, each of which divides into sub-stages. The pseudometanauplius stage is exclusive to species that lay their eggs within an ovigerous sac: so-called "sac-spawners". The larvae grow and moult repeatedly as they develop, replacing their rigid exoskeleton when it becomes too small. Smaller animals moult more frequently than larger ones. Yolk reserves within their body nourish the larvae through metanauplius stage.
By the calyptopsis stages differentiation has progressed far enough for them to develop a mouth and a digestive tract, and they begin to eat phytoplankton. By that time their yolk reserves are exhausted and the larvae must have reached the photic zone, the upper layers of the ocean where algae flourish. During the furcilia stages, segments with pairs of swimmerets are added, beginning at the frontmost segments. Each new pair becomes functional only at the next moult. The number of segments added during any one of the furcilia stages may vary even within one species depending on environmental conditions. After the final furcilia stage, an immature juvenile emerges in a shape similar to an adult, and subsequently develops gonads and matures sexually.
Reproduction
During the mating season, which varies by species and climate, the male deposits a sperm sack at the female's genital opening (named thelycum). The females can carry several thousand eggs in their ovary, which may then account for as much as one third of the animal's body mass. Krill can have multiple broods in one season, with interbrood intervals lasting on the order of days.
Krill employ two types of spawning mechanism. The 57 species of the genera Bentheuphausia, Euphausia, Meganyctiphanes, Thysanoessa, and Thysanopoda are "broadcast spawners": the female releases the fertilised eggs into the water, where they usually sink, disperse, and are on their own. These species generally hatch in the nauplius 1 stage, but have recently been discovered to hatch sometimes as metanauplius or even as calyptopis stages. The remaining 29 species of the other genera are "sac spawners", where the female carries the eggs with her, attached to the rearmost pairs of thoracopods until they hatch as metanauplii, although some species like Nematoscelis difficilis may hatch as nauplius or pseudometanauplius.
Moulting
Moulting occurs whenever a specimen outgrows its rigid exoskeleton. Young animals, growing faster, moult more often than older and larger ones. The frequency of moulting varies widely by species and is, even within one species, subject to many external factors such as latitude, water temperature, and food availability. The subtropical species Nyctiphanes simplex, for instance, has an overall inter-moult period of two to seven days: larvae moult on the average every four days, while juveniles and adults do so, on average, every six days. For E. superba in the Antarctic sea, inter-moult periods ranging between 9 and 28 days depending on the temperature between have been observed, and for Meganyctiphanes norvegica in the North Sea the inter-moult periods range also from 9 and 28 days but at temperatures between . E. superba is able to reduce its body size when there is not enough food available, moulting also when its exoskeleton becomes too large. Similar shrinkage has also been observed for E. pacifica, a species occurring in the Pacific Ocean from polar to temperate zones, as an adaptation to abnormally high water temperatures. Shrinkage has been postulated for other temperate-zone species of krill as well.
Lifespan
Some high-latitude species of krill can live for more than six years (e.g., Euphausia superba); others, such as the mid-latitude species Euphausia pacifica, live for only two years. Subtropical or tropical species' longevity is still shorter, e.g., Nyctiphanes simplex, which usually lives for only six to eight months.
Swarming
Most krill are swarming animals; the sizes and densities of such swarms vary by species and region. For Euphausia superba, swarms reach 10,000 to 60,000 individuals per cubic metre. Swarming is a defensive mechanism, confusing smaller predators that would like to pick out individuals. In 2012, Gandomi and Alavi presented what appears to be a successful stochastic algorithm for modelling the behaviour of krill swarms. The algorithm is based on three main factors: " (i) movement induced by the presence of other individuals (ii) foraging activity, and (iii) random diffusion."
Vertical migration
Krill typically follow a diurnal vertical migration. It has been assumed that they spend the day at greater depths and rise during the night toward the surface. The deeper they go, the more they reduce their activity, apparently to reduce encounters with predators and to conserve energy. Swimming activity in krill varies with stomach fullness. Sated animals that had been feeding at the surface swim less actively and therefore sink below the mixed layer. As they sink they produce feces which employs a role in the Antarctic carbon cycle. Krill with empty stomachs swim more actively and thus head towards the surface.
Vertical migration may be a 2–3 times daily occurrence. Some species (e.g., Euphausia superba, E. pacifica, E. hanseni, Pseudeuphausia latifrons, and Thysanoessa spinifera) form surface swarms during the day for feeding and reproductive purposes even though such behaviour is dangerous because it makes them extremely vulnerable to predators.
Experimental studies using Artemia salina as a model suggest that the vertical migrations of krill several hundreds of metres, in groups tens of metres deep, could collectively create enough downward jets of water to have a significant effect on ocean mixing.
Dense swarms can elicit a feeding frenzy among fish, birds and mammal predators, especially near the surface. When disturbed, a swarm scatters, and some individuals have even been observed to moult instantly, leaving the exuvia behind as a decoy.
Krill normally swim at a pace of 5–10 cm/s (2–3 body lengths per second), using their swimmerets for propulsion. Their larger migrations are subject to ocean currents. When in danger, they show an escape reaction called lobstering—flicking their caudal structures, the telson and the uropods, they move backwards through the water relatively quickly, achieving speeds in the range of 10 to 27 body lengths per second, which for large krill such as E. superba means around . Their swimming performance has led many researchers to classify adult krill as micro-nektonic life-forms, i.e., small animals capable of individual motion against (weak) currents. Larval forms of krill are generally considered zooplankton.
Biogeochemical cycles
The Antarctic krill is an important species in the context of biogeochemical cycling and in the Antarctic food web. It plays a prominent role in the Southern Ocean because of its ability to cycle nutrients and to feed penguins and baleen and blue whales.
Human uses
Harvesting history
Krill have been harvested as a food source for humans and domesticated animals since at least the 19th century, and possibly earlier in Japan, where it was known as okiami. Large-scale fishing developed in the late 1960s and early 1970s, and now occurs only in Antarctic waters and in the seas around Japan. Historically, the largest krill fishery nations were Japan and the Soviet Union, or, after the latter's dissolution, Russia and Ukraine. The harvest peaked, which in 1983 was about 528,000 tonnes in the Southern Ocean alone (of which the Soviet Union took in 93%), is now managed as a precaution against overfishing.
In 1993, two events caused a decline in krill fishing: Russia exited the industry; and the Convention for the Conservation of Antarctic Marine Living Resources (CCAMLR) defined maximum catch quotas for a sustainable exploitation of Antarctic krill. After an October 2011 review, the Commission decided not to change the quota.
The annual Antarctic catch stabilised at around 100,000 tonnes, which is roughly one fiftieth of the CCAMLR catch quota. The main limiting factor was probably high costs along with political and legal issues. The Japanese fishery saturated at some 70,000 tonnes.
Although krill are found worldwide, fishing in Southern Oceans are preferred because the krill are more "catchable" and abundant in these regions. Particularly in Antarctic seas which are considered as pristine, they are considered a "clean product".
In 2018 it was announced that almost every krill fishing company operating in Antarctica will abandon operations in huge areas around the Antarctic Peninsula from 2020, including "buffer zones" around breeding colonies of penguins.
Human consumption
Although the total biomass of Antarctic krill may be as abundant as 400 million tonnes, the human impact on this keystone species is growing, with a 39% increase in total fishing yield to 294,000 tonnes over 2010–2014. Major countries involved in krill harvesting are Norway (56% of total catch in 2014), the Republic of Korea (19%), and China (18%).
Krill is a rich source of protein and omega-3 fatty acids which are under development in the early 21st century as human food, dietary supplements as oil capsules, livestock food, and pet food. Krill tastes salty with a somewhat stronger fish flavor than shrimp. For mass consumption and commercially prepared products, they must be peeled to remove the inedible exoskeleton.
In 2011, the US Food and Drug Administration published a letter of no objection for a manufactured krill oil product to be generally recognized as safe (GRAS) for human consumption.
Krill (and other planktonic shrimp, notably Acetes spp.) are most widely consumed in Southeast Asia, where it is fermented (with the shells intact) and usually ground finely to make shrimp paste. It can be stir-fried and eaten paired with white rice or used to add umami flavors to a wide variety of traditional dishes. The liquid from the fermentation process is also harvested as fish sauce.
Bio-inspired robotics
Krill are agile swimmers in the intermediate Reynolds number regime, in which there are not many solutions for uncrewed underwater robotics, and have inspired robotic platforms to both study their locomotion as well as find design solutions for underwater robots.
| Biology and health sciences | Crustaceans | null |
231386 | https://en.wikipedia.org/wiki/Navel | Navel | The navel (clinically known as the umbilicus; : umbilici or umbilicuses; commonly known as the belly button or tummy button) is a protruding, flat, or hollowed area on the abdomen at the attachment site of the umbilical cord.
Structure
The umbilicus is used to visually separate the abdomen into quadrants.
The umbilicus is a prominent scar on the abdomen, with its position being relatively consistent among humans. The skin around the waist at the level of the umbilicus is supplied by the tenth thoracic spinal nerve (T10 dermatome). The umbilicus itself typically lies at a vertical level corresponding to the junction between the L3 and L4 vertebrae, with a normal variation among people between the L3 and L5 vertebrae.
Parts of the adult navel include the "umbilical cord remnant" or "umbilical tip", which is the often protruding scar left by the detachment of the umbilical cord. This is located in the center of the navel, sometimes described as the belly button. Around the cord remnant is the "umbilical collar", formed by the dense fibrous umbilical ring. Surrounding the umbilical collar is the periumbilical skin. Directly behind the navel is a thick fibrous cord formed from the umbilical cord, called the urachus, which originates from the bladder.
The belly button is unique to each individual due to it being a scar, and various general forms have been classified by medical practitioners.
Outie: A navel consisting of the umbilical tip protruding past the periumbilical skin is an outie. Essentially any navel which is not concave.
Swirly/spiral: A rare form in which the umbilical cord scar forms a swirl shape.
Split: The protruding umbilical cord scar extends outwards, but is cleft in two by a fissure which extends part or all the way through the umbilical cord scar. This form is similar in appearance to a coffee bean.
Protrusion: The umbilical cord remnant is completely divulged, exposing the full umbilical scar.
Circlet: Although the entirety of the umbilical cord remnant sits out with the umbilical collar, the centre of the knot is inset by a deep fissure. Unlike a split outie, in this form the fissure is contained centrally and does not extend past the umbilical cord remnant in any direction, much akin to a 'donut' shape.
Innie: A navel in which the umbilical tip does not protrude past the periumbilical skin. Any navel which is concave.
Round: Round navels are completely circular with no hooding.
Vertical: Some navels present in the form of a more elongate hollow parallel with the linea alba.
Oval: This form consists of three variants; superior hooding, inferior hooding, no hooding.
T-shaped: As the name states, the scar is in the shape of a T, and may have superior hooding to various extent.
Horizontal: The scar is the least visible, as the natural lines of the tendinous intersection fold over the scar.
Distorted: Any navel which does not fit well into any of the other categories.
Clinical significance
Disorders
Outies are sometimes mistaken for umbilical hernias; however, they are a completely different shape with no health concern, unlike an umbilical hernia. The navel (specifically abdominal wall) would be considered an umbilical hernia if the protrusion were 5 centimeters or more. The diameter of an umbilical hernia is usually 1/2-inch or more. While the shape of the human navel may be affected by long term changes to diet and exercise, unexpected change in shape may be the result of ascites.
In addition to change in shape being a possible side effect from ascites and umbilical hernias, the navel can be involved in umbilical sinus or fistula, which in rare cases can lead to menstrual or fecal discharge from the navel. Menstrual discharge from the umbilicus is a rare disorder associated with umbilical endometriosis.
Other disorders
Omphalitis, an inflammatory condition of the umbilicus in the newborn, usually caused by a bacterial infection.
Omphalophobia is the fear of belly buttons. People suffering from omphalophobia are terrified of belly buttons—their own or, in some cases, those of others. They do not like touching their belly buttons (or even other people touching it). Sometimes just seeing a belly button is enough to make them feel disgusted or terrified.
Surgery
To minimize scarring, the navel is a recommended site of incision for various surgeries, including transgastric appendicectomy, gall bladder surgery, and the umbilicoplasty procedure itself.
Other animals
All placental mammals have a navel, although it is generally more conspicuous in humans.
Fashion, society and culture
The public exposure of the male and female midriff and bare navel was considered taboo at times in the past in Western cultures, being considered immodest or indecent. Female navel exposure was banned in some jurisdictions, but community perceptions have changed to this now being acceptable. The crop top is a shirt that often exposes the belly button and has become more common among young people. Exposure of the male navel has rarely been stigmatised and has become particularly popular in recent years, due to the strong resurgence of the male crop top and male navel piercing. The navel and midriff are often also displayed in bikinis, or when low-rise pants are worn.
While the West was relatively resistant to navel-baring clothing until the 1980s, it has long been a fashion with Indian women, often displayed with saris or lehengas.
The Japanese have long had a special regard for the navel. During the early Jōmon period in northern Japan, three small balls indicating the breasts and navel were pasted onto flat clay objects to represent the female body. The navel was exaggerated in size, informed by the belief that the navel symbolized the center where life began.
In Arabic-Levantine culture, belly dancing is a popular art form that consists of dance movements focused on the torso and navel.
Buddhism and Hinduism refer to the chakra of the navel as the manipura. In qigong, the navel is seen as the main energy centre, or dantian. In Hinduism, the Kundalini energy is sometimes described as being located at the navel.
| Biology and health sciences | External anatomy and regions of the body | Biology |
19290728 | https://en.wikipedia.org/wiki/Perovskite | Perovskite | Perovskite (pronunciation: ) is a calcium titanium oxide mineral composed of calcium titanate (chemical formula ). Its name is also applied to the class of compounds which have the same type of crystal structure as , known as the perovskite structure, which has a general chemical formula . Many different cations can be embedded in this structure, allowing the development of diverse engineered materials.
History
The mineral was discovered in the Ural Mountains of Russia by Gustav Rose in 1839 and is named after Russian mineralogist Lev Perovski (1792–1856). Perovskite's notable crystal structure was first described by Victor Goldschmidt in 1926 in his work on tolerance factors. The crystal structure was later published in 1945 from X-ray diffraction data on barium titanate by Helen Dick Megaw.
Occurrence
Found in the Earth's mantle, perovskite's occurrence at Khibina Massif is restricted to the silica under-saturated ultramafic rocks and foidolites, due to the instability in a paragenesis with feldspar. Perovskite occurs as small anhedral to subhedral crystals filling interstices between the rock-forming silicates.
Perovskite is found in contact carbonate skarns at Magnet Cove, Arkansas, US, in altered blocks of limestone ejected from Mount Vesuvius, in chlorite and talc schist in the Urals and Switzerland, and as an accessory mineral in alkaline and mafic igneous rocks, nepheline syenite, melilitite, kimberlites and rare carbonatites. Perovskite is a common mineral in the Ca-Al-rich inclusions found in some chondritic meteorites.
The stability of perovskite in igneous rocks is limited by its reaction relation with sphene. In volcanic rocks perovskite and sphene are not found together, the only exception being an etindite from Cameroon.
A rare-earth-bearing variety knopite with the chemical formula is found in alkali intrusive rocks in the Kola Peninsula and near Alnö, Sweden. A niobium-bearing variety dysanalyte occurs in carbonatite near Schelingen, Kaiserstuhl, Germany.
In stars and brown dwarfs
In stars and brown dwarfs the formation of perovskite grains is responsible for the depletion of titanium oxide in the photosphere. Stars with a low temperature have dominant bands of TiO in their spectrum; as the temperature gets lower for stars and brown dwarfs with an even lower mass, forms and at temperatures below 2000 K TiO is undetectable. The presence of TiO is used to define the transition between cool M-dwarf stars and the colder L-dwarfs.
Physical properties
The eponymous Perovskite crystallizes in the Pbnm space group (No. 62) with lattice constants a = 5.39 Å, b = 5.45 Å and c = 7.65 Å.
Perovskites have a nearly cubic structure with the general formula . In this structure the A-site ion, in the center of the lattice, is usually an alkaline earth or rare-earth element. B-site ions, on the corners of the lattice, are 3d, 4d, and 5d transition metal elements. The A-site cations are in 12-fold coordination with the anions, while the B-site cations are in 6-fold coordination. A large number of metallic elements are stable in the perovskite structure if the Goldschmidt tolerance factor t is in the range of 0.75 to 1.0.
where RA, RB and RO are the ionic radii of A and B site elements and oxygen, respectively. The stability of perovskites can be characterized with the tolerance and octahedral factors. When conditions are not fulfilled, a layered geometry for edge-sharing or face-sharing octahedra or lower B-site coordination is preferred. These are good structural bounds, but not an empirical prediction.
Perovskites have sub-metallic to metallic luster, colorless streak, and cube-like structure along with imperfect cleavage and brittle tenacity. Colors include black, brown, gray, orange to yellow. Perovskite crystals may appear to have the cubic crystal form, but are often pseudocubic and actually crystallize in the orthorhombic system, as is the case for (strontium titanate, with the larger strontium cation in the A-site, is cubic). Perovskite crystals have been mistaken for galena; however, galena has a better metallic luster, greater density, perfect cleavage and true cubic symmetry.
Perovskite derivatives
Double perovskites
A double perovskite has a formula of and replaces half the B sites with B, where A are alkaline or rare earth metals and B are transition metals. The cation arrangement will differ based on charge, coordination geometry, and the ratio between A cation and B cation radii. The B and B cations lead to different ordering schemes. These ordering schemes are rock salt, columnar, and layered structures.
Rock salt is an alternating, three-dimensional checkerboard of B and B' polyhedra. This structure is the most common from an electrostatic point of view, as the B sites will have different valence states. Columnar arrangement can be viewed as sheets of B-cation polyhedral viewed from the [111] direction. Layered structures are seen as sheets of B and B polyhedra.
Lower dimensional perovskites
3D perovskites form when there is a smaller cation in the A site so octahedra can be corner shared. 2D perovskites form when the A-site cation is larger so octahedra sheets are formed. In 1D perovskites, a chain of octahedra is formed
while in 0D perovskites, individual octahedra are separated from each other. Both 1D and 0D perovskites lead to quantum confinement and are investigated for lead-free perovskite solar cell materials.
| Physical sciences | Minerals | Earth science |
19294750 | https://en.wikipedia.org/wiki/Phylum | Phylum | In biology, a phylum (; : phyla) is a level of classification or taxonomic rank below kingdom and above class. Traditionally, in botany the term division has been used instead of phylum, although the International Code of Nomenclature for algae, fungi, and plants accepts the terms as equivalent. Depending on definitions, the animal kingdom Animalia contains about 31 phyla, the plant kingdom Plantae contains about 14 phyla, and the fungus kingdom Fungi contains about eight phyla. Current research in phylogenetics is uncovering the relationships among phyla within larger clades like Ecdysozoa and Embryophyta.
General description
The term phylum was coined in 1866 by Ernst Haeckel from the Greek (, "race, stock"), related to (, "tribe, clan"). Haeckel noted that species constantly evolved into new species that seemed to retain few consistent features among themselves and therefore few features that distinguished them as a group ("a self-contained unity"): "perhaps such a real and completely self-contained unity is the aggregate of all species which have gradually evolved from one and the same common original form, as, for example, all vertebrates. We name this aggregate [a] [i.e., stock / tribe] ()." In plant taxonomy, August W. Eichler (1883) classified plants into five groups named divisions, a term that remains in use today for groups of plants, algae and fungi.
The definitions of zoological phyla have changed from their origins in the six Linnaean classes and the four of Georges Cuvier.
Informally, phyla can be thought of as groupings of organisms based on general specialization of body plan. At its most basic, a phylum can be defined in two ways: as a group of organisms with a certain degree of morphological or developmental similarity (the phenetic definition), or a group of organisms with a certain degree of evolutionary relatedness (the phylogenetic definition). Attempting to define a level of the Linnean hierarchy without referring to (evolutionary) relatedness is unsatisfactory, but a phenetic definition is useful when addressing questions of a morphological nature—such as how successful different body plans were.
Definition based on genetic relation
The most important objective measure in the above definitions is the "certain degree" that defines how different organisms need to be members of different phyla. The minimal requirement is that all organisms in a phylum should be clearly more closely related to one another than to any other group. Even this is problematic because the requirement depends on knowledge of organisms' relationships: as more data become available, particularly from molecular studies, we are better able to determine the relationships between groups. So phyla can be merged or split if it becomes apparent that they are related to one another or not. For example, the bearded worms were described as a new phylum (the Pogonophora) in the middle of the 20th century, but molecular work almost half a century later found them to be a group of annelids, so the phyla were merged (the bearded worms are now an annelid family). On the other hand, the highly parasitic phylum Mesozoa was divided into two phyla (Orthonectida and Rhombozoa) when it was discovered the Orthonectida are probably deuterostomes and the Rhombozoa protostomes.
This changeability of phyla has led some biologists to call for the concept of a phylum to be abandoned in favour of placing taxa in clades without any formal ranking of group size.
Definition based on body plan
A definition of a phylum based on body plan has been proposed by paleontologists Graham Budd and Sören Jensen (as Haeckel had done a century earlier). The definition was posited because extinct organisms are hardest to classify: they can be offshoots that diverged from a phylum's line before the characters that define the modern phylum were all acquired. By Budd and Jensen's definition, a phylum is defined by a set of characters shared by all its living representatives.
This approach brings some small problems—for instance, ancestral characters common to most members of a phylum may have been lost by some members. Also, this definition is based on an arbitrary point of time: the present. However, as it is character based, it is easy to apply to the fossil record. A greater problem is that it relies on a subjective decision about which groups of organisms should be considered as phyla.
The approach is useful because it makes it easy to classify extinct organisms as "stem groups" to the phyla with which they bear the most resemblance, based only on the taxonomically important similarities. However, proving that a fossil belongs to the crown group of a phylum is difficult, as it must display a character unique to a sub-set of the crown group. Furthermore, organisms in the stem group of a phylum can possess the "body plan" of the phylum without all the characteristics necessary to fall within it. This weakens the idea that each of the phyla represents a distinct body plan.
A classification using this definition may be strongly affected by the chance survival of rare groups, which can make a phylum much more diverse than it would be otherwise.
Known phyla
Animals
Total numbers are estimates; figures from different authors vary wildly, not least because some are based on described species, some on extrapolations to numbers of undescribed species. For instance, around 25,000–27,000 species of nematodes have been described, while published estimates of the total number of nematode species include 10,000–20,000; 500,000; 10 million; and 100 million.
Plants
The kingdom Plantae is defined in various ways by different biologists (see Current definitions of Plantae). All definitions include the living embryophytes (land plants), to which may be added the two green algae divisions, Chlorophyta and Charophyta, to form the clade Viridiplantae. The table below follows the influential (though contentious) Cavalier-Smith system in equating "Plantae" with Archaeplastida, a group containing Viridiplantae and the algal Rhodophyta and Glaucophyta divisions.
The definition and classification of plants at the division level also varies from source to source, and has changed progressively in recent years. Thus some sources place horsetails in division Arthrophyta and ferns in division Monilophyta, while others place them both in Monilophyta, as shown below. The division Pinophyta may be used for all gymnosperms (i.e. including cycads, ginkgos and gnetophytes), or for conifers alone as below.
Since the first publication of the APG system in 1998, which proposed a classification of angiosperms up to the level of orders, many sources have preferred to treat ranks higher than orders as informal clades. Where formal ranks have been provided, the traditional divisions listed below have been reduced to a very much lower level, e.g. subclasses.
Fungi
Phylum Microsporidia is generally included in kingdom Fungi, though its exact relations remain uncertain, and it is considered a protozoan by the International Society of Protistologists (see Protista, below). Molecular analysis of Zygomycota has found it to be polyphyletic (its members do not share an immediate ancestor), which is considered undesirable by many biologists. Accordingly, there is a proposal to abolish the Zygomycota phylum. Its members would be divided between phylum Glomeromycota and four new subphyla incertae sedis (of uncertain placement): Entomophthoromycotina, Kickxellomycotina, Mucoromycotina, and Zoopagomycotina.
Protists
Kingdom Protista (or Protoctista) is included in the traditional five- or six-kingdom model, where it can be defined as containing all eukaryotes that are not plants, animals, or fungi. Protista is a paraphyletic taxon, which is less acceptable to present-day biologists than in the past. Proposals have been made to divide it among several new kingdoms, such as Protozoa and Chromista in the Cavalier-Smith system.
Protist taxonomy has long been unstable, with different approaches and definitions resulting in many competing classification schemes. Many of the phyla listed below are used by the Catalogue of Life, and correspond to the Protozoa-Chromista scheme, with updates from the latest (2022) publication by Cavalier-Smith. Other phyla are used commonly by other authors, and are adapted from the system used by the International Society of Protistologists (ISP). Some of the descriptions are based on the 2019 revision of eukaryotes by the ISP.
The number of protist phyla varies greatly from one classification to the next. The Catalogue of Life includes Rhodophyta and Glaucophyta in kingdom Plantae, but other systems consider these phyla part of Protista. In addition, less popular classification schemes unite Ochrophyta and Pseudofungi under one phylum, Gyrista, and all alveolates except ciliates in one phylum Myzozoa, later lowered in rank and included in a paraphyletic phylum Miozoa. Even within a phylum, other phylum-level ranks appear, such as the case of Bacillariophyta (diatoms) within Ochrophyta. These differences became irrelevant after the adoption of a cladistic approach by the ISP, where taxonomic ranks are excluded from the classifications after being considered superfluous and unstable. Many authors prefer this usage, which lead to the Chromista-Protozoa scheme becoming obsolete.
Bacteria
Currently there are 41 bacterial phyla (not including "Cyanobacteria") that have been validly published according to the Bacteriological Code
Abditibacteriota
Acidobacteriota, phenotypically diverse and mostly uncultured
Actinomycetota, High-G+C Gram positive species
Aquificota, deep-branching
Armatimonadota
Atribacterota
Bacillota, Low-G+C Gram positive species, such as the spore-formers Bacilli (aerobic) and Clostridia (anaerobic)
Bacteroidota
Balneolota
Bdellovibrionota
Caldisericota, formerly candidate division OP5, Caldisericum exile is the sole representative
Calditrichota
Campylobacterota
Chlamydiota
Chlorobiota, green sulphur bacteria
Chloroflexota, green non-sulphur bacteria
Chrysiogenota, only 3 genera (Chrysiogenes arsenatis, Desulfurispira natronophila, Desulfurispirillum alkaliphilum)
Coprothermobacterota
Deferribacterota
Deinococcota, Deinococcus radiodurans and Thermus aquaticus are "commonly known" species of this phyla
Dictyoglomota
Elusimicrobiota, formerly candidate division Thermite Group 1
Fibrobacterota
Fusobacteriota
Gemmatimonadota
Ignavibacteriota
Kiritimatiellota
Lentisphaerota, formerly clade VadinBE97
Mycoplasmatota, notable genus: Mycoplasma
Myxococcota
Nitrospinota
Nitrospirota
Planctomycetota
Pseudomonadota, the most well-known phylum, containing species such as Escherichia coli or Pseudomonas aeruginosa
Rhodothermota
Spirochaetota, species include Borrelia burgdorferi, which causes Lyme disease
Synergistota
Thermodesulfobacteriota
Thermomicrobiota
Thermotogota, deep-branching
Verrucomicrobiota
Archaea
Currently there are 2 phyla that have been validly published according to the Bacteriological Code
Nitrososphaerota
Thermoproteota, second most common archaeal phylum
Other phyla that have been proposed, but not validly named, include:
"Euryarchaeota", most common archaeal phylum
"Korarchaeota"
"Nanoarchaeota", ultra-small symbiotes, single known species
| Biology and health sciences | Taxonomic rank | Biology |
12561015 | https://en.wikipedia.org/wiki/4chan | 4chan | 4chan is an anonymous English-language imageboard website. Launched by Christopher "moot" Poole in October 2003, the site hosts boards dedicated to a wide variety of topics, from video games and television to literature, cooking, weapons, music, history, technology, anime, physical fitness, politics, and sports, among others. Registration is not available, except for staff, and users typically post anonymously. , 4chan receives more than 22 million unique monthly visitors, of whom approximately half are from the United States.
4chan was created as an unofficial English-language counterpart to the Japanese imageboard Futaba Channel, also known as 2chan, and its first boards were originally used for posting images and discussion related to anime. The site has been described as a hub of Internet subculture, its community being influential in the formation and popularization of prominent Internet memes, such as lolcats, Rickrolling, rage comics, wojaks, Pepe the Frog, as well as hacktivist and political movements, such as Anonymous and the alt-right.
4chan has often been the subject of media attention as a source of controversies, including the coordination of pranks and harassment against websites and Internet users, and the posting of illegal and offensive content as a result of its lax censorship and moderation policies. In 2008, The Guardian summarized the 4chan community as "lunatic, juvenile [...] brilliant, ridiculous and alarming".
Background
The majority of posting on 4chan takes place on imageboards, on which users have the ability to share images and create threaded discussions. , the site's homepage lists 75 imageboards and one Flash animation board. Most boards have their own set of rules and are dedicated to a specific topic, including anime and manga, video games, music, literature, fitness, politics, and sports, among others. Uniquely, the "Random" board—also known as /b/—enforces few rules.
4chan is the Internet's most trafficked imageboard, according to the Los Angeles Times. 4chan's Alexa rank was 853 in March 2022, though it has been as high as 56. It is provided to its users free of charge and consumes a large amount of bandwidth; as a result, its financing has often been problematic. Poole has acknowledged that donations alone could not keep the site online, and turned to advertising to help make ends meet. However, the explicit content hosted on 4chan has deterred businesses who do not want to be associated with the site's content. In January 2009, Poole signed a new deal with an advertising company; in February 2009, he was $20,000 in debt, and the site was continuing to lose money. The 4chan servers were moved from Texas to California in August 2008, which upgraded the maximum bandwidth throughput of 4chan from 100 Mbit/s to 1 Gbit/s.
Unlike most web forums, 4chan does not have a registration system, allowing users to post anonymously. Posting is ephemeral, as threads receiving recent replies are "bumped" to the top of their respective board and old threads are deleted as new ones are created. Any nickname may be used when posting, even one that has been previously adopted, such as "Anonymous" or "moot". In place of registration, 4chan has provided tripcodes as an optional form of authenticating a poster's identity. As making a post without filling in the "Name" field causes posts to be attributed to "Anonymous", general understanding on 4chan holds that Anonymous is not a single person but a collective (hive) of users.
Moderators generally post without a name even when performing sysop actions. A "capcode" may be used to attribute the post to "Anonymous ## Mod", although moderators often post without the capcode. In a 2011 interview on Nico Nico Douga, Poole explained that there are approximately 20 volunteer moderators active on 4chan. 4chan also has a junior moderation team, called "janitors", who may delete posts or images and suggest that the normal moderation team ban a user, but who cannot post with a capcode. Revealing oneself as a janitor is grounds for immediate dismissal. Gianluca Stringhini, an associate professor at Boston University College of Engineering, said in August 2024, "The only moderation on the platform appears to be for clearly illegal content, such as child pornography. Everything else remains untouched."
4chan has been the target of occasional denial of service attacks. For instance, on December 28, 2010, 4chan and other websites went down due to such an attack, following which Poole said on his blog, "We now join the ranks of MasterCard, Visa, PayPal, etc and is an exclusive club!"
History
The site was launched as 4chan.net on October 1, 2003, by Christopher Poole, a then-15-year-old student from New York City using the online handle "moot". Poole had been a regular participant on Something Awful's subforum "Anime Death Tentacle Rape Whorehouse" (ADTRW), where many users were familiar with the Japanese imageboard format and Futaba Channel ("2chan.net"). When creating 4chan, Poole obtained Futaba Channel's open source code and translated the Japanese text into English using AltaVista's Babel Fish online translator. After the site's creation, Poole invited users from the ADTRW subforum, many of whom were dissatisfied with the site's moderation, to visit 4chan, which he advertised as an English-language counterpart to Futaba Channel and a place for Western fans to discuss anime and manga. At its founding, the site only hosted one board: /b/ (Anime/Random).
Before the end of 2003, several new anime-related boards were added, including /h/ (Hentai), /c/ (Anime/Cute), /d/ (Hentai/Alternative), /w/ (Wallpapers/Anime), /y/ (Yaoi), and /a/ (Anime). In the early days of the website, Poole hosted meetings from 2005 to 2008 in various locations to promote it, such as Otakon, that popularized some of the first 4chan-related memes.
Additionally, a lolicon board was created at /l/ (Lolikon), but was disabled following the posting of genuine child pornography and ultimately deleted in October 2004, after threats of legal action. In February 2004, GoDaddy suspended the 4chan.net domain, prompting Poole to move the site to its current domain at 4chan.org. On March 1, 2004, Poole announced that he lacked the funds to pay the month's server bill, but was able to continue operations after receiving a swarm of donations from users. In June 2004, 4chan experienced six weeks of downtime due to PayPal suspending 4chan's donations service after receiving complaints about the site's content.
Following 4chan's return, several non-anime related boards were introduced, including /k/ (Weapons), /o/ (Auto), and /v/ (Video Games). In 2008, nine new boards were created, including the sports board at /sp/, the fashion board at /fa/ and the "Japan/General" (the name later changed to "Otaku Culture") board at /jp/. By this point, 4chan's culture had altered, moving away from the "early, more childish," humour, as evident by the likes of Project Chanology; trolling underwent a so-called "golden age" that took aim at American corporate media.
In January 2011, Poole announced the deletion of the /r9k/ ("ROBOT9000") and /new/ (News) boards, saying that /new/ had become devoted to racist discussions, and /r9k/ no longer served its original purpose of being a test implementation of xkcd's ROBOT9000 script. During the same year, the /soc/ board was created in an effort to reduce the number of socialization threads on /b/. /r9k/ was restored on October 23, 2011, along with /hc/ ("Hardcore", previously deleted), /pol/ (a rebranding of /new/) and the new /diy/ board, in addition to an apology by Poole where he recalls how he criticized the deletion of Encyclopedia Dramatica and realized that he had done the same.
In 2010, 4chan had implemented reCAPTCHA in an effort to thwart spam arising from JavaScript worms. By November 2011, 4chan made the transition to utilizing Cloudflare following a series of DDoS attacks. The 4chan imageboards were rewritten in valid HTML5/CSS3 in May 2012 in an effort to improve client-side performance. On September 28, 2012, 4chan introduced a "4chan pass" that, when purchased, "allows users to bypass typing a reCAPTCHA verification when posting and reporting posts on the 4chan image boards"; the money raised from the passes to go towards supporting the site.
On January 21, 2015, Poole stepped down as the site's administrator, citing stress from controversies such as Gamergate as the reason for his departure. On September 21, 2015, Poole announced that Hiroyuki Nishimura had purchased from him the ownership rights to 4chan, without disclosing the terms of the acquisition. Nishimura was the former administrator of 2channel between 1999 and 2014, the website forming the basis for anonymous posting culture which influenced later websites such as Futaba Channel and 4chan; Nishimura lost 2channel's domain after it was seized by his registrar, Jim Watkins due the latter's alleged financial difficulties. Wired later reported that Japanese toy manufacturer Good Smile Company, Japanese telecommunication Dwango, and Nishimura's company Future Search Brazil may have helped facilitate Nishimura's purchase, with anonymous sources telling the publication that Good Smile obtained partial ownership in the website as compensation.
In October 2016, it was reported that the site was facing financial difficulties that could lead to its closure or radical changes. In a post titled "Winter is Coming", Hiroyuki Nishimura explained, "We had tried to keep 4chan as is. But I failed. I am sincerely sorry", citing server costs, infrastructure costs, and network fees.
On November 17, 2018, it was announced that the site would be split into two, with the work-safe boards moved to a new domain, 4channel.org, while the NSFW boards would remain on the 4chan.org domain. In a series of posts on the topic, Nishimura explained that the split was due to 4chan being blacklisted by most advertising companies and that the new 4channel domain would allow for the site to receive advertisements by mainstream ad providers. All boards returned to the 4chan.org domain in December 2023 for unknown reasons, and 4channel.org now redirects to 4chan.org.
In a 2020 interview with Vice Media, several current or past moderators spoke about what they perceived as racist intent behind the site's management. They alleged that a managing moderator named RapeApe was attempting to use the site as a recruitment tool for the alt-right, and that Nishimura was "hands-off, leaving moderation of the site primarily to RapeApe." Neither Nishimura nor RapeApe responded to these allegations. Far-right extremism has been reported by public authorities, commentators and civil society groups as connected, in part, to 4chan, an association that had arisen by 2015. According to 4chan's filings to the New York Attorney General's Office, 4chan signed an agreement to pay RapeApe $3,000 a month for their services in 2015. By May 2022, that fee had risen to $4,400 a month. The submitted documents also revealed RapeApe lamenting that 4chan was "getting the shaft" over the Buffalo terrorist attack and his attempt to persuade the advertising platform Bid.Glass to reverse their exit from the website.
Christopher Poole
Poole concealed his real-life identity until it was revealed on July 9, 2008, in The Wall Street Journal. Prior to that, he had used the alias "moot".
In April 2009, an open Internet poll conducted by Time magazine voted Poole as the world's most influential person of 2008. The results were questioned even before the poll completed, as automated voting programs and manual ballot stuffing were used to influence the vote. 4chan's interference with the vote seemed increasingly likely, when it was found that reading the first letter of the first 21 candidates in the poll spelled out a phrase containing two 4chan memes: "mARBLECAKE. ALSO, THE GAME."
On September 12, 2009, Poole gave a talk regarding 4chan's reputation as a "Meme Factory" at the Paraflows Symposium in Vienna, Austria, which was part of the Paraflows 09 festival, themed Urban Hacking. In this talk, Poole mainly attributed this both to the anonymous system and to the lack of data retention on the site ("The site has no memory.").
In April 2010, Poole testified in the trial United States of America v. David Kernell as a government witness, explaining the terminology used on 4chan to the prosecutor, ranging from "OP" to "lurker", as well as the nature of the data given to the FBI as part of the search warrant, including how users can be uniquely identified from site audit logs.
Notable boards
/b/
The "random" board, /b/, follows the design of Futaba Channel's Nijiura ("Random") board. It was the first board created, and has been described as 4chan's most popular board, accounting for 30% of site traffic in 2009. Gawker's Nick Douglas summarized /b/ as a board where "people try to shock, entertain, and coax free porn from each other." /b/ has a "no rules" policy, except for bans on certain illegal content, such as child pornography, invasions of other websites (posting floods of disruptive content), and under-18 viewing, all of which are inherited from site-wide rules. The "no invasions" rule was added in late 2006, after /b/ users spent most of that summer "invading" Habbo Hotel. The "no rules" policy also applies to actions of administrators and moderators, which means that users may be banned at any time, for any reason, including for no reason at all. Due partially to its anonymous nature, board moderation is not always successful—indeed, the site's anti-child pornography rule is a subject of jokes on /b/. Christopher Poole told The New York Times, in a discussion on the moderation of /b/, that "the power lies in the community to dictate its own standards" and that site staff simply provided a framework.
The humor of /b/'s many users, who refer to themselves as "/b/tards", is often incomprehensible to newcomers and outsiders, and is characterized by intricate inside jokes and dark comedy. Users often refer to each other, and much of the outside world, as fags. They are often referred to by outsiders as trolls, who regularly act with the intention of "doing it for the lulz", a corruption of "LOL" used to denote amusement at another's expense. A significant amount of media coverage is in response to /b/'s culture, which has been characterized as adolescent, crude and spiteful, with one publication writing that their "bad behavior is encouraged by the site's total anonymity and the absence of an archive". Douglas cited Encyclopedia Dramatica's definition of /b/ as "the asshole of the Internets [sic]". Mattathias Schwartz of The New York Times likened /b/ to "a high-school bathroom stall, or an obscene telephone party line", while Baltimore City Paper wrote that "in the high school of the Internet, /b/ is the kid with a collection of butterfly knives and a locker full of porn." Wired describes /b/ as "notorious".
Each post is assigned a post number. Certain post numbers are sought after with a large amount of posting taking place to "GET" them. A "GET" occurs when a post's number ends in a special number, such as 12345678, 22222222, or every millionth post. A sign of 4chan's scaling, according to Poole, was when GETs lost meaning due to the high post rate resulting in a GET occurring every few weeks. He estimated /b/'s post rate in July 2008 to be 150,000–200,000 posts per day.
/mu/
The music board, /mu/, is dedicated to the discussion of music artists, albums, genres, and instruments. Described as "4chan's best kept secret" and a "surprisingly artistic side of 4chan", /mu/ is used by users to share their music interests with similar minds and discover "great music they would never have found otherwise" with many moments of insightful candor that can affirm or challenge their own musical tastes. The board has gained notoriety for earnestly focusing upon and promoting challenging and otherwise obscure music. Some common genres discussed on /mu/ include shoegaze, experimental hip hop, witch house, IDM, midwest emo, vaporwave, and K-pop. There is a significant overlap between user bases of /mu/ and music site Rate Your Music. The board's culture has inspired many online music communities and meme pages on social media that emulate /mu/'s posting style.
Publications such as Pitchfork and Entertainment Weekly noted the board played a significant role in popularizing various music artists, such as Death Grips, Neutral Milk Hotel, Car Seat Headrest, and Have a Nice Life. Prominent music critic Anthony Fantano began his career on /mu/ and developed a significant following there. Some artists, like Zeal & Ardor and Conrad Tao, admitted to posting their music anonymously on /mu/ to get honest feedback, as well as find inspiration from the board. In particular, Zeal & Ardor said their sound, which mixes black metal with spirituals, came from suggestions by two users. Andrew W.K. did a Q&A with the board's users in 2011, causing the servers to crash from the increased traffic. Death Grips seeded various clues on /mu/ in 2012 about their then-upcoming albums The Money Store and No Love Deep Web. A rendition of "Royals" by Lorde appeared on /mu/ in 2012 before its official release, although she denied ever writing on the board in 2014. Singer Lauren Mayberry shared on Twitter in 2015 a link to a thread on /mu/ about her band's song "Leave a Trace" to showcase what online misogyny looks like. An alleged unreleased Radiohead song, titled "Putting Ketchup in the Fridge" and "How Do You Sit Still", was initially reported as genuine by NME and Spin until CNN revealed it was a hoax promoted by the board's users.
The board has been acknowledged for sharing rare music recordings and unreleased materials, as well as finding albums thought to be lost. Notable examples include the works of Duster, D>E>A>T>H>M>E>T>A>L by Panchiko, and All Lights Fucked on the Hairy Amp Drooling by Godspeed You! Black Emperor. This was described by NPR as resembling "a secret club of preservationists obsessed with the articulation of a near-dead language". The board has attracted further attention for various projects done by its users. A group called The Pablo Collective posted a 4-track remix album of Kanye West's The Life of Pablo titled The Death of Pablo to /mu/, claiming it was based on a recurring dream from one of the board's users. A role-playing game based on Neutral Milk Hotel's In the Aeroplane Over the Sea, designed with help from the board's users, received coverage from Polygon and Pitchfork.
/pol/
/pol/ ("Politically Incorrect") is 4chan's political discussion board. A stickied thread on its front page states that the board's intended purpose is "discussion of news, world events, political issues, and other related topics." /pol/ was created in October 2011 as a rebranding of 4chan's news board, /new/, which was deleted that January for a high volume of racist discussion.
Although there had previously been a strong left-libertarian contingent to 4chan activists, there was a gradual rightward turn on 4chan's politics board in the early-mid 2010s, with the fundamentalist approach to free speech contributing. The board quickly attracted posters with a political persuasion that later would be described with a new term, the alt-right. Media sources have characterized /pol/ as predominantly racist and sexist, with many of its posts taking an explicitly neo-Nazi bent. The site's far-reaching culture of vitriolic and discriminatory content is "most closely associated" with /pol/, although only it features predominant Alt-Right beliefs; /pol/, like other boards, has been prominent in the dissemination of memes, in cases, featuring coordination to disperse Alt-Right sentiments. /pol/ "increasingly became synonymous with 4chan as a whole". The Southern Poverty Law Center regards /pol/'s rhetorical style as widely emulated by white supremacist websites such as The Daily Stormer; the Stormers editor, Andrew Anglin, concurred. /pol/ was where screenshots of Trayvon Martin's hacked social media accounts were initially posted. The board's users have started antifeminist, homophobic, transphobic, and anti-Arab Twitter campaigns.
Many /pol/ users favored Donald Trump during his 2016 United States presidential campaign. Both Trump and his son, Donald Trump Jr., appeared to acknowledge the support by tweeting /pol/-associated memes. Upon his successful election, a /pol/ moderator embedded a pro-Trump video at the top of all of the board's pages.
/r9k/
/r9k/ is a board that implements Randall Munroe's "ROBOT9000" algorithm, where no exact reposts are permitted. It is credited as the origin of the "greentext" rhetorical style which often center around stories of social interactions and resulting ineptness. By 2012, personal confession stories of self-loathing, depression, and attempted suicide began to supersede /b/-style roleplaying, otaku, and video game discussion.
It became a popular gathering place for the controversial online incel community. The "beta uprising" or "beta rebellion" meme, the idea of taking revenge against women, jocks and others perceived as the cause of incels' problems, was popularized on the subsection. The perpetrator of the Toronto van attack referenced 4chan and an incel rebellion in a Facebook post he made prior to the attack, while praising self-identified incel Elliot Rodger, the killer behind the 2014 Isla Vista killings. He claims to have talked with both Harper-Mercer and Rodger on Reddit and 4chan and believes that he was part of a "beta uprising", also posting a message on 4chan about his intention the day before his attack.
/sci/
/sci/ is 4chan's science and mathematics board. On September 26, 2011, an anonymous user on /sci/ posted a question regarding the shortest possible way to watch all possible orders of episodes of the anime The Melancholy of Haruhi Suzumiya in nonchronological order. Shortly after, an anonymous user responded with a mathematical proof that argued viewers would have to watch at least 93,884,313,611 episodes to see all possible orderings. Seven years later, professional mathematicians recognized the mathematical proof as a partial solution to a superpermutations problem that was unsolved for 25 years. Australian mathematician Greg Egan later published a proof inspired by the proof from the anonymous 4chan user, both of which are recognized as significant advances to the problem.
/v/
/v/ is 4chan's video games board. The board has spawned multiple Internet memes, most notably the NPC Wojak in 2016 (derived from the gaming term non-player character to describe those who do not think for themselves or make their own conscious decisions).
/x/
The "paranormal" board, /x/, is dedicated to discussing topics regarding unexplained phenomena, the supernatural, and non-political conspiracy theories. /x/ was initially launched in January 2005 as 4chan's general photography board; in February 2007, it was repurposed as a paranormal-themed board.
Many of the earliest creepypastas (Internet horror-related legends) were created on /x/. The idea of the Backrooms gained popularity thanks to a thread on /x/ created on 12 May 2019, where the users were asked to "post disquieting images that just feel 'off'." There, the first photo depicting the Backrooms was uploaded and another user commented on it with the first story about the Backrooms, claiming that one enters the Backrooms when they "noclip out of reality in the wrong areas". After the 4chan post gained fame, several Internet users wrote horror stories relating to the Backrooms. Many memes were created and shared across social media, further popularizing the creepypasta.
American model Allison Harvard first gained notoriety in 2005 as an Internet meme on the /x/ board where she became known as Creepy Chan. Known for her large eyes and peculiar interests like fascination with blood, photos she posted on her blog were widely circulated on the board. She gained mainstream notoriety in 2009 and again in 2011 by appearing on America's Next Top Model. She would visit /x/ after new episodes of America's Next Top Model would air to see what was being written about her and participate in the discussions.
The SCP Foundation, a fictional secret organization documented by the collaborative writing wiki project of the same name, originated on /x/ in 2007, when the first SCP file, SCP-173, was posted by an anonymous user. Initially a stand-alone short story, many additional SCP files were created shortly after; these new SCPs copied SCP-173's style and were set within the same fictional universe. A stand-alone wiki was created in January 2008 on the EditThis wiki hosting service to display the SCP articles. The EditThis website did not have moderators, or the ability to delete articles. Members communicated through individual article talk pages and the /x/ board.
/x/ was the first place where the 2015 viral video 11B-X-1371 was posted. The board also contributed to investigating and popularizing the controversial Sad Satan video game.
Internet culture
Early internet memes
"[A] significant and influential element of contemporary internet culture", 4chan is responsible for many early memes and the site has received positive attention for its association with memes. This included "So I herd u liek mudkipz", which involved a phrase based on Pokémon and which generated numerous YouTube tribute videos, and the term "an hero" as a synonym for suicide, after a misspelling in the Myspace online memorial of seventh grader Mitchell Henderson. 4chan and other websites, such as the satirical Encyclopedia Dramatica, have also contributed to the development of significant amounts of leetspeak.
A lolcat is an image combining a photograph of a cat with solecistic text intended to contribute humour, widely popularized by 4chan in the form of a weekly post dedicated to them and a corresponding theme.
In 2005, the installment of a word filter which changed "egg" to "duck", and thus "eggroll" to "duckroll", across 4chan led to a bait-and-switch meme in which users deceitfully linked to a picture of a duck on wheels. This was then modified into users linking to the music video for Rick Astley's 1987 song "Never Gonna Give You Up". Thus, the "rickroll" was born.
A link to the YouTube video of Tay Zonday's song "Chocolate Rain" was posted on /b/ on July 11, 2007, and then subsequently circulated by users, becoming a very popular internet meme. The portion of the song in which Zonday turns away from the microphone, with a caption stating "I move away from the mic to breathe in", became an oft-repeated meme on 4chan and inspired remixes. Fellow YouTuber Boxxy's popularity was also due in part to 4chan.
In his American incarnation, Pedobear is an anthropomorphic bear child predator that is often used within the community to mock contributors showing a sexual interest in children. Pedobear is one of the most popular memes on non-English imageboards, and has gained recognition across Europe, appearing in offline publications. It has been used as a symbol of pedophilia by Maltese graffiti vandals prior to a papal visit.
Anonymous and anti-Scientology activism
4chan has been labeled as the starting point of the Anonymous meme by The Baltimore City Paper, due to the norm of posts signed with the "Anonymous" moniker. The National Posts David George-Cosh said it has been "widely reported" that Anonymous is associated with 4chan and 711chan, as well as numerous Internet Relay Chat (IRC) channels.
Through its association with Anonymous, 4chan has become associated with Project Chanology, a worldwide protest against the Church of Scientology held by members of Anonymous. On January 15, 2008, a 4chan user posted to /b/, suggesting participants "do something big" against the Church of Scientology's website. This message resulted in the Church receiving threatening phone calls. It quickly grew into a large real-world protest. Unlike previous Anonymous attacks, this action was characterized by 4chan memes including rickrolls and Guy Fawkes masks. The raid drew criticism from some 4chan users who felt it would bring the site undesirable attention.
My Little Pony: Friendship is Magic fandom
The adult fandom and subculture dedicated to the children's animated television series My Little Pony: Friendship Is Magic began on the "Comics & Cartoons" (/co/) board of 4chan. The show was first discussed with some interest around its debut in October 2010. The users of /co/ took a heightened interest in the show after a critical Cartoon Brew article was shared, resulting in praise for its plot, characters, and animation style. Discussion of the show extended to /b/, eventually to a point of contention. Discussion then spread forth to communities external to 4chan, including the establishment of the fan websites, causing the show to reach a wider audience across the internet.
"This post is art"
On July 30, 2014, an anonymous user made a reply in a thread on the board /pol/ "Politically Incorrect" of 4chan, criticizing modern art in an ironic fashion, saying:
Less than an hour later the post was photographed off the screen and framed by another user who posted another reply in the thread with a photo of the framed quote. Later the user, after endorsement by other anonymous users in the thread, created an auction on eBay for the framed photo which quickly rose to high prices, culminating in a price of $90,900.
Controversies and harassment incidents
Internet raids
According to The Washington Post, "the site's users have managed to pull off some of the highest-profile collective actions in the history of the Internet."
Users of 4chan and other websites "raided" Hal Turner by launching DDoS attacks and prank calling his phone-in radio show during December 2006 and January 2007. The attacks caused Turner's website to go offline. This cost thousands of dollars of bandwidth bills according to Turner. In response, Turner sued 4chan, 7chan, and other websites; however, he lost his plea for an injunction and failed to receive letters from the court.
KTTV Fox 11 aired a report on Anonymous, calling them a group of "hackers on steroids", "domestic terrorists", and collectively an "Internet hate machine" on July 26, 2007. Slashdot founder Rob Malda posted a comment made by another Slashdot user, Miang, stating that the story focused mainly on users of "4chan, 7chan and 420chan". Miang claimed that the report "seems to confuse /b/ raids and motivational poster templates with a genuine threat to the American public", arguing that the "unrelated" footage of a van exploding shown in the report was to "equate anonymous posting with domestic terror".
On July 10, 2008, the swastika CJK unicode character (卐) appeared at the top of Google's Hot Trends list—a tally of the most used search terms in the United States—for several hours. It was later reported that the HTML numeric character reference for the symbol had been posted on /b/, with a request to perform a Google search for the string. A multitude of /b/ visitors followed the order and pushed the symbol to the top of the chart, though Google later removed the result.
Later that year, the private Yahoo! Mail account of Sarah Palin, Republican vice presidential candidate in the 2008 United States presidential election, was hacked by a 4chan user. The hacker posted the account's password on /b/, and screenshots from within the account to WikiLeaks. A /b/ user then logged in and changed the password, posting a screenshot of him sending an email to a friend of Palin's informing her of the new password on the /b/ thread. However, he forgot to blank out the password in the screenshot. A multitude of /b/ users attempted to log in with the new password, and the account was automatically locked out by Yahoo!. The incident was criticized by some /b/ users. One user commented, "seriously, /b/. We could have changed history and failed, epically." The FBI and Secret Service began investigating the incident shortly after its occurrence. On September 20 it was revealed they were questioning David Kernell, the son of Democratic Tennessee State Representative Mike Kernell.
The stock price of Apple Inc. fell significantly in October 2008 after a hoax story was submitted to CNN's user-generated news site iReport.com claiming that company CEO Steve Jobs had suffered a major heart attack. The source of the story was traced back to 4chan.
In May 2009, members of the site attacked YouTube, posting pornographic videos on the video-sharing platform under names of teenage celebrities. The attack spawned the popular Internet meme and catchphrase "I'm 12 years old and what is this?" as a response to a user comment on one such video. A 4chan member acknowledged being part of the attack, telling BBC News that it was in response to YouTube "deleting music". In January 2010, members of the site attacked YouTube again in response to the suspension of YouTube user lukeywes1234 for failing to meet the minimum age requirement of thirteen. The videos uploaded by the user had apparently become popular with 4chan members, who subsequently became angered after the account was suspended and called for a new wave of pornographic videos to be uploaded to YouTube on January 6, 2010. Later the same year, 4chan made numerous disruptive pranks directed at singer Justin Bieber.
In September 2010, in retaliation against the Bollywood film industry's hiring of Aiplex Software to launch cyberattacks against The Pirate Bay, Anonymous members, recruited through posts on 4chan boards, subsequently initiated their own attacks, dubbed Operation Payback, targeting the website of the Motion Picture Association of America and the Recording Industry Association of America. The targeted websites usually went offline for a short period of time due to the attacks, before recovering.
The website of the UK law firm ACS:Law, which was associated with an anti-piracy client, was affected by the cyber-attack. In retaliation for the initial attacks being called only a minor nuisance, Anonymous launched more attacks, bringing the site down yet again. After coming back up, the front page accidentally revealed a backup file of the entire website, which contained over 300 megabytes of private company emails, which were leaked to several torrents and across several sites on the Internet. It was suggested that the data leak could cost the law firm up to £500,000 in fines for breaching British Data Protection Laws.
In January 2011, BBC News reported that the law firm announced they were to stop "chasing illegal file-sharers". Head of ACS:Law Andrew Crossley in a statement to a court addressed issues which influenced the decision to back down "I have ceased my work ... I have been subject to criminal attack. My e-mails have been hacked. I have had death threats and bomb threats."
In August 2012, 4chan users attacked a third-party sponsored Mountain Dew campaign, Dub the Dew, where users were asked to submit and vote on name ideas for a green apple flavor of the drink. Users submitted entries such as "Diabeetus", "Fapple", several variations of "Gushing Granny", and "Hitler did nothing wrong".
Threats of violence
On October 18, 2006, the Department of Homeland Security warned National Football League officials in Miami, New York City, Atlanta, Seattle, Houston, Oakland, and Cleveland about a possible threat involving the simultaneous use of dirty bombs at stadiums. The threat claimed that the attack would be carried out on October 22, the final day of the Muslim holy month of Ramadan. Both the FBI and the Department of Homeland Security expressed doubt concerning the credibility of the threats, but warned the relevant organizations as a precaution. The threat turned out to be an ill-conceived hoax perpetrated by a grocery store clerk in Wisconsin with no terrorist ties. The FBI considered it a clearly frivolous threat and the 20-year-old man was charged with fabricating a terrorist threat, sentenced to six months in prison followed by six months' house arrest, and ordered to pay $26,750 in restitution.
Around midnight on September 11, 2007, a student posted photographs of mock pipe bombs and another photograph of him holding them while saying he would blow up his high school—Pflugerville High School in Pflugerville, Texas—at 9:11 am on September 11. Users of 4chan helped to track him down by finding the perpetrator's father's name in the Exif data of a photograph he took, and contacted the police. He was arrested before school began that day. The incident turned out to be a hoax; the "weapons" were toys and there were no actual bombs.
A 20-year-old from Melbourne, Australia, was arrested on December 8, 2007, after apparently posting on 4chan that he was "going to shoot and kill as many people as I can until which time I am incapacitated or killed by the police". The post, accompanied by an image of another man holding a shotgun, threatened a shopping mall near Beverly Hills. While the investigation was still open, he was charged with criminal defamation for a separate incident but died before the case was heard.
On February 4, 2009, a posting on the 4chan /b/ board said there would be a school shooting at St Eskils Gymnasium in Eskilstuna, Sweden, leading 1,250 students and 50 teachers to be evacuated. A 21-year-old man was arrested after 4chan provided the police with the IP address of the poster. Police said that the suspect called it off as a joke, and they released him after they found no indication that the threat was serious.
On June 28, 2018, a man was arrested following an indictment by the U.S. Department of Justice "on one count of transmitting in interstate and foreign commerce a threat to injure the person of another." The indictment alleged that he posted anonymously to /pol/ the day after the Unite the Right rally, communicating an intention to attack protestors at an upcoming right-wing demonstration, ostensibly to elicit sympathy for the alt-right movement. "I'm going to bring a Remington 700 and start shooting Alt-right guys. We need sympathy after that landwhale got all the liberals teary eyed, so someone is going to have to make it look like the left is becoming more violent and radicalized. It's a false flag for sure, but I'll be aiming for the more tanned/dark haired muddied jeans in the crowd so real whites won't have to worry," he wrote, according to the indictment.
In 2023, a 38-year-old of Monmouth Junction, New Jersey, was arrested for threatening Volusia County, Florida sheriff Mike Chitwood on 4chan due to Chitwood's condemnation of anti-Semitism. According to authorities, the poster, who lived 974 miles away from Volusia County, advocated "shoot[ing] Chitwood in the head and murder[ing] him" in a February 22 post.
In April of that same year, two other 4chan users, residents of California and Connecticut respectively, were also arrested for threatening to kill Chitwood on 4chan.
Celebrity photo leaks
On August 31, 2014, a compromise of user passwords at iCloud allowed a large number of private photographs taken by celebrities to be posted online, initially on 4chan. As a result of the incident, 4chan announced that it would enforce a Digital Millennium Copyright Act policy, which would allow content owners to remove material that had been shared on the site illegally, and would ban users who repeatedly posted stolen material.
Gamergate
Also in August 2014, 4chan was involved in the Gamergate controversy, which began with unsubstantiated allegations about indie game developer Zoë Quinn from an ex-boyfriend, followed by false allegations from anonymous Internet users. The allegations were followed by a harassment campaign against several women in the video game industry, organized by 4chan users, particularly /r9k/. Discussion regarding Gamergate was banned on 4chan due to alleged rule violations, and Gamergate supporters moved to alternate forums such as 8chan.
Murder in Port Orchard, Washington
According to court documents filed on November 5, 2014, there were images posted to 4chan that appeared to be of a murder victim. The body was discovered in Port Orchard, Washington, after the images were posted. The posts were accompanied by the text: "Turns out it's way harder to strangle someone to death than it looks on the movies." A later post said: "Check the news for Port Orchard, Washington, in a few hours. Her son will be home from school soon. He'll find her, then call the cops. I just wanted to share the pics before they find me." The victim was Amber Lynn Coplin, aged 30. The suspect, 33-year-old David Michael Kalac, surrendered to police in Oregon later the same day; he was charged with second-degree murder involving domestic violence. Kalac was convicted in April 2017 and was sentenced to 82 years in prison the following month.
Death of Jeffrey Epstein
A report of Jeffrey Epstein's death was posted on /pol/ around 40 minutes before ABC News broke the news. It was originally suspected that the unidentified person who made the posts may have been a first responder, prompting a review by the New York City Fire Department, who later stated that the post did not come from a member of its department.
2022 Buffalo shooting
On May 14, 2022, a mass shooting occurred at a supermarket in Buffalo, New York, US. The accused, Payton S. Gendron, is reported to have written a racist manifesto released May 12 (two days before the shooting), with the manifesto including birth date and other biographical details, that match the suspect in custody. The author wrote that he began to frequent 4chan, including its Politically Incorrect message board /pol/, beginning in May 2020, where he was exposed to the Great Replacement conspiracy theory.
ISP bans
AT&T temporary ban
On July 26, 2009, AT&T's DSL branch temporarily blocked access to the img.4chan.org domain (host of /b/ and /r9k/), which was initially believed to be an attempt at Internet censorship, and met with hostility on 4chan's part. The next day, AT&T issued a statement claiming that the block was put in place after an AT&T customer was affected by a DoS attack originating from IP addresses connected to img.4chan.org, and was an attempt to "prevent this attack from disrupting service for the impacted AT&T customer, and... our other customers." AT&T maintains that the block was not related to the content on 4chan.
4chan's founder Christopher Poole responded with the following:
Major news outlets have reported that the issue may be related to the DDoS-ing of 4chan, and that 4chan users suspected the then-owner of Swedish-based website Anontalk.com.
Verizon temporary ban
On February 4, 2010, 4chan started receiving reports from Verizon Wireless customers that they were having difficulties accessing the site's image boards. After investigating, Poole found out that only the traffic on port 80 to the boards.4chan.org domain was affected, leading members to believe that the block was intentional. Three days later, Verizon Wireless confirmed that 4chan was "explicitly blocked". The block was lifted several days later.
Telstra ban
On March 20, 2019, Australian telecom company Telstra denied access to millions of Australians to 4chan, 8chan, Zero Hedge and LiveLeak as a reaction to the Christchurch mosque shootings.
New Zealand
Following the Christchurch mosque shootings, numerous ISPs temporarily blocked any site hosting a copy of the livestream of the shooting, including 4chan. The ISPs included Spark, Vodafone, Vocus and 2degrees.
| Technology | Social network and blogging | null |
4582654 | https://en.wikipedia.org/wiki/Passenger | Passenger | A passenger is a person who travels in a vehicle, but does not bear any responsibility for the tasks required for that vehicle to arrive at its destination or otherwise operate the vehicle, and is not a steward. The vehicles may be bicycles, buses, cars, passenger trains, airliners, ships, ferryboats, personal watercraft, all terrain vehicles, snowmobiles, and other methods of transportation.
Crew members (if any), as well as the driver or pilot of the vehicle, are usually not considered to be passengers. For example, a flight attendant on an airline would not be considered a passenger while on duty and the same with those working in the kitchen or restaurant on board a ship as well as cleaning staff, but an employee riding in a company car being driven by another person would be considered a passenger, even if the car was being driven on company business.
Legal status
In most jurisdictions, laws have been enacted that dictate the legal obligations of the owner of a vehicle or vessel, or of the driver or pilot of the same, towards the passengers. With respect to passengers riding in cars and vans, the driver may owe a duty of care to passengers, particularly where the passenger's presence in the vehicle can be seen to "confer some benefit on the driver other than the benefit of his or her company or the mere sharing of expenses". In other situations, however, guest statutes may limit the ability of passengers to sue the driver of the vehicle over an accident. Many places require cars to be outfitted with measures specifically for the protection of passengers, such as passenger-side air bags. With respect to passengers on commercial vehicles or vessels, both national laws and international treaties require that the carrier act with a certain standard of care. The number of passengers that a vehicle or vessel may legally carry is defined as its seating capacity.
Terms
Revenue passenger
A revenue passenger is someone who has paid a transport operator for her or his trip. That excludes non-paying passengers such as airline employees flying on free or nearly-free passes, babies and children who do not have a seat of their own, etc. However, passengers who paid for their trip with a frequent-flyer program mileage award are usually included.
This term is used in the transportation industry, in particular in traffic measures such as revenue passenger kilometer (RPK) and revenue passenger mile (RPM).
Revenue passenger kilometres
Revenue passenger kilometres (RPKs) and revenue passenger miles (RPMs) are measures of traffic for an airline flight, bus, or train calculated by multiplying the number of revenue-paying passengers aboard the vehicle by the distance traveled. On long-distance buses and trains (and some planes), passengers may board and disembark at intermediate stops, in which case RPMs/RPKs have to be calculated for each segment if a careful total is needed.
Revenue passenger miles can be considered the basic amount of "production" that an airline creates. The revenue passenger miles can be compared to the available seat miles over an airline's system to determine the overall passenger load factor. These measurements can further be used to measure unit revenues and unit costs.
No pax
In transportation, a "no pax" trip is a trip without passengers. For example, no-pax flights are Air cargo, ferry and positioning flights. Similarly, with a public transit bus it can be used at the beginning and end of a driver’s work shift to/from the bus terminal, or in the non-commute leg of a commuter bus service.
In such cases, the main display signs on the front and curbside of the bus typically display a message such as “no pax” or “out of service” (sometimes abbreviated as “O/S”).
British railway passenger train categories
In British railway parlance, passenger, as well as being the end user of a service, is also a categorisation of the type of rolling stock used. In the British case, there are several categories of passenger train, which include:
Express passenger, which constitutes long distance and high speed railway travel between major locations such as ports and cities.
Semi-fast express passenger, a type of service that is high speed, though stops at selected destinations of high population density en route.
Local passenger, the lowest category of British passenger train, which provides a service that stops at all stations between major destinations, for the benefit of local populations.
| Technology | Basics_7 | null |
4585070 | https://en.wikipedia.org/wiki/Non-cellular%20life | Non-cellular life | Non-cellular life, also known as acellular life, is life that exists without a cellular structure for at least part of its life cycle. Historically, most definitions of life postulated that an organism must be composed of one or more cells, but, for some, this is no longer considered necessary, and modern criteria allow for forms of life based on other structural arrangements.
Nucleic acid-containing infectious agents
Viruses
Researchers initially described viruses as "poisons" or "toxins", then as "infectious proteins"; but they possess genetic material, a defined structure, and the ability to spontaneously assemble from their constituent parts. This has spurred extensive debate as to whether they should be regarded as fundamentally organic or inorganic — as very small biological organisms or as very large biochemical molecules. Without their hosts, they are not able to perform any of the functions of life - such as respiration, growth, or reproduction. Since the 1950s, many scientists have thought of viruses as existing at the border between chemistry and life; a gray area between living and nonliving.
Viroids
If viruses are borderline cases or nonliving, viroids are further from being living organisms. Viroids are some of the smallest infectious agents, consisting solely of short strands of circular, single-stranded RNA without protein coats. They are only known to infect flowering plants, of which some are of commercial importance. Viroid genomes are extremely small in size, ranging from 246 to 467 nucleobases. In comparison, the genome of the smallest viruses capable of causing an infection are around 2,000 nucleobases in size. Viroid RNA does not code for any protein. Its replication mechanism hijacks RNA polymerase II, a host-cell enzyme normally associated with synthesis of messenger RNA from DNA, which instead catalyzes "rolling circle" synthesis of new RNA using the viroid's RNA as a template. Some viroids are ribozymes, having catalytic properties which allow self-cleavage and ligation of unit-size genomes from larger replication intermediates.
A possible explanation of the origin of viroids sees them as "living relics" from a hypothetical, ancient, and non-cellular RNA world before the evolution of DNA or of protein. This view, first proposed in the 1980s, regained popularity in the 2010s to explain crucial intermediate steps in the evolution of life from inanimate matter (abiogenesis).
Obelisks
In 2024, researchers announced the possible discovery of viroid-like, but distinct, RNA-based elements dubbed obelisks. Obelisks, found in sequence databases of the human microbiome, are possibly hosted in gut bacteria. They differ from viroids in that they code for two distinct proteins, dubbed "oblins", and for the predicted rod-like secondary structure of their RNA.
First universal common ancestor
The first universal common ancestor is an example of a proposed non-cellular lifeform, as it is the earliest ancestor of the last universal common ancestor, its sister lineages, and every currently living cell.
| Biology and health sciences | Biology basics | Biology |
4586930 | https://en.wikipedia.org/wiki/Lineation%20%28geology%29 | Lineation (geology) | Lineations in structural geology are linear structural features within rocks. There are several types of lineations, intersection lineations, crenulation lineations, mineral lineations and stretching lineations being the most common. Lineation field measurements are recorded as map lines with a plunge angle and azimuth.
Intersection lineations
Intersection lineations are linear structures formed by the intersection of any two surfaces in a three-dimensional space. The trace of bedding on an intersecting foliation plane commonly appears as colour stripes generally parallel to local fold's hinges. Intersection lineations can also be due to the intersection of two foliations.
Intersection lineations are measured in relation to the two structures which intersect to form them. For instance, according to the measurement conventions of structural geology, original bedding, S0 intersected by a fold's axial plane foliation, forms an intersection lineation L0-1, with an azimuth and plunge defined by the fold. This is the typical cleavage-bedding intersection angle and is diagnostic of the plunge of the fold on all parts of the fold.
Stretching lineations
Stretching lineations are formed by shearing of rocks during asymmetric deformation of a rock mass. Stretching lineations record primarily the vector of greatest stretch, which is perpendicular to the principle plane of shortening.
A stretching lineation may be visualised as a ball of treacle (molasses) which, when pulled, forms a cigar-shaped rod parallel to the direction in which it is pulled. This is parallel to the direction in which a shearing force, as found in a shear zone, stretches the rock. Shortening occurs at the same time as elongation but in a perpendicular sense to the stretched rod.
With reference to the image at right (top), the conglomerate pebbles most likely were deposited as sub-spherical pebbles and boulders. During deformation the rock was flattened and then stretched by movement along a ductile shear zone within which this outcrop resides. The spherical conglomerate pebbles stretched along the direction of movement of this shear zone, attaining their current somewhat flattened cigar-shaped form. The pebbles thus record important information on the orientation of the shear zone (subvertical) and the direction of movement of the shear zone, and the overall change in pebble shape from originally sub-spherical to presently elongate cigar-shaped, allows one to quantify the strain experienced by the rock mass in the geologic past.
Stretching lineations may also manifest as linear features upon pre-existing surfaces such as foliations within shear zones (see image at right, below). In such a case the lineation may not be as obvious in plan and may require measurement as a rake upon a planar surface. In this case, the two lineations are formed in the same deformation event but are manifest differently owing to the different rheologies of the deformed rocks.
Finally, the key difference between a stretching lineation and an intersection lineation is that stretching lineations carry no information on the orientation of other planar fabrics within a rock mass. In the case of the illustrated lineations within the sandstone, they do not record an earlier deformation event's foliation and cannot be used to infer orientation information for folds or original bedding.
Linear structures are extremely important in structural mapping, they can be used to separate deformation phases and to determine the kinematics of deformation. Quartz rods are one of the most eye-catching linear structures in deformed rocks. Despite relatively rare, rods are described in several places worldwide. Wherever rods occur, they are promptly noticed. Rods form a conspicuous coarse lineation, frequently highly contrasting with the surrounding rock in regions that was under high strain. The term rod or rodding, in geology, broadly refers to a mass of rock, which has assumed a cylindrical shape while accommodating strain; however, different definitions are found in the literature. The mechanisms of rod formation can be constrained from field observations. They are frequently parallel to fold axes and lie at right angles to the direction of maximum compression. For more information on RODS please refer to:
| Physical sciences | Structural geology | Earth science |
100340 | https://en.wikipedia.org/wiki/Millipede | Millipede | Millipedes (originating from the Latin , "thousand", and , "foot") are a group of arthropods that are characterised by having two pairs of jointed legs on most body segments; they are known scientifically as the class Diplopoda, the name derived from this feature. Each double-legged segment is a result of two single segments fused together. Most millipedes have very elongated cylindrical or flattened bodies with more than 20 segments, while pill millipedes are shorter and can roll into a tight ball. Although the name "millipede" derives from Latin for "thousand feet", no species was known to have 1,000 or more until the discovery in 2020 of Eumillipes persephone, which can have over 1,300 legs. There are approximately 12,000 named species classified into 16 orders and around 140 families, making Diplopoda the largest class of myriapods, an arthropod group which also includes centipedes and other multi-legged creatures.
Most millipedes are slow-moving detritivores, eating decaying leaves and other dead plant matter; however, some eat fungi or drink plant fluid. Millipedes are generally harmless to humans, although some can become household or garden pests. Millipedes can be an unwanted nuisance particularly in greenhouses where they can potentially cause severe damage to emergent seedlings. Most millipedes defend themselves with a variety of chemicals secreted from pores along the body, although the tiny bristle millipedes are covered with tufts of detachable bristles. Its primary defence mechanism is to curl into a tight coil, thereby protecting its legs and other vital delicate areas on the body behind a hard exoskeleton. Reproduction in most species is carried out by modified male legs called gonopods, which transfer packets of sperm to females.
First appearing in the Silurian period, millipedes are some of the oldest known land animals. Some members of prehistoric groups, such as Arthropleura, grew to over ; the largest modern species reach maximum lengths of . The longest extant species is the giant African millipede (Archispirostreptus gigas).
Among myriapods, millipedes have traditionally been considered most closely related to the tiny pauropods, although some molecular studies challenge this relationship. Millipedes can be distinguished from the somewhat similar but only distantly related centipedes (class Chilopoda), which move rapidly, are venomous, carnivorous, and have only a single pair of legs on each body segment.
The scientific study of millipedes is known as diplopodology, and a scientist who studies them is called a diplopodologist.
Etymology and names
The term "millipede" is widespread in popular and scientific literature, but among North American scientists, the term "milliped" (without the terminal e) is also used. Other vernacular names include "thousand-legger" or simply "diplopod". The science of millipede biology and taxonomy is called diplopodology: the study of diplopods.
Classification
Approximately 12,000 millipede species have been described. Estimates of the true number of species on earth range from 15,000 to as high as 80,000. Few species of millipede are at all widespread; they have very poor dispersal abilities, depending as they do on terrestrial locomotion and humid habitats. These factors have favoured genetic isolation and rapid speciation, producing many lineages with restricted ranges.
The living members of the Diplopoda are divided into sixteen orders in two subclasses. The basal subclass Penicillata contains a single order, Polyxenida (bristle millipedes). All other millipedes belong to the subclass Chilognatha consisting of two infraclasses: Pentazonia, containing the short-bodied pill millipedes, and Helminthomorpha (worm-like millipedes), containing the great majority of the species.
Outline of classification
The higher-level classification of millipedes is presented below, based on Shear, 2011, and Shear & Edgecombe, 2010 (extinct groups). Recent cladistic and molecular studies have challenged the traditional classification schemes above, and in particular the position of the orders Siphoniulida and Polyzoniida is not yet well established. The placement and positions of extinct groups (†) known only from fossils is tentative and not fully resolved. After each name is listed the author citation: the name of the person who coined the name or defined the group, even if not at the current rank.
Class Diplopoda de Blainville in Gervais, 1844
Subclass Penicillata Latreille, 1831
Order Polyxenida Verhoeff, 1934
Subclass †Arthropleuridea (placed in Penicillata by some authors)
Order †Arthropleurida Waterlot, 1934
Order †Eoarthropleurida Shear & Selden, 1995
Order †Microdecemplicida Wilson & Shear, 2000
Subclass Chilognatha Latreille, 1802
Order †Zosterogrammida Wilson, 2005 (Chilognatha incertae sedis)
Infraclass Pentazonia Brandt, 1833
Order †Amynilyspedida Hoffman, 1969
Superorder Limacomorpha Pocock, 1894
Order Glomeridesmida Cook, 1895
Superorder Oniscomorpha Pocock, 1887
Order Glomerida Brandt, 1833
Order Sphaerotheriida Brandt, 1833
Infraclass Helminthomorpha Pocock, 1887
Superorder †Archipolypoda Scudder, 1882
Order †Archidesmida Wilson & Anderson 2004
Order †Cowiedesmida Wilson & Anderson 2004
Order †Euphoberiida Hoffman, 1969
Order †Palaeosomatida Hannibal & Krzeminski, 2005
Order †Pleurojulida Schneider & Werneburg, 1998 (possibly sister to Colobognatha)
Subterclass Colobognatha Brandt, 1834
Order Platydesmida Cook, 1895
Order Polyzoniida Cook, 1895
Order Siphonocryptida Cook, 1895
Order Siphonophorida Newport, 1844
Subterclass Eugnatha Attems, 1898
Superorder Juliformia Attems, 1926
Order Julida Brandt, 1833
Order Spirobolida Cook, 1895
Order Spirostreptida Brandt, 1833
Superfamily †Xyloiuloidea Cook, 1895 (Sometimes aligned with Spirobolida)
Superorder Nematophora Verhoeff, 1913
Order Callipodida Pocock, 1894
Order Chordeumatida Pocock 1894
Order Stemmiulida Cook, 1895
Order Siphoniulida Cook, 1895
Superorder Merocheta Cook, 1895
Order Polydesmida Pocock, 1887
Evolution
Millipedes are among the first animals to have colonised land during the Silurian period. Early forms probably ate mosses and primitive vascular plants. There are two major groups of millipedes whose members are all extinct: the Archipolypoda ("ancient, many-legged ones") which contain the oldest known terrestrial animals, and Arthropleuridea, which contain the largest known land invertebrates. Pneumodesmus newmani is the earliest member of the millipedes from the late Wenlock epoch of the late Silurian around , or early Lochkovian of the early Devonian around 414 million years ago, known from long fragment and has clear evidence of spiracles (breathing holes) attesting to its air-breathing habits. Other early fossils of millipedes are Kampecaris obanensis and Archidesmus sp. from 425 millions years ago in the late Silurian. During the Carboniferous, Arthropleura became the largest known land-dwelling invertebrate on record, length exceeding . Reason of gigantism of Arthropleura is not clearly known, previously considered that is due to high oxygen levels, but later studies consider that is more likely because of lack of competition. Millipedes also exhibit the earliest evidence of chemical defence, as some Devonian fossils have defensive gland openings called ozopores.
Living groups
The history of scientific millipede classification began with Carl Linnaeus, who in his 10th edition of Systema Naturae, 1758, named seven species of Julus as "Insecta Aptera" (wingless insects). In 1802, the French zoologist Pierre André Latreille proposed the name Chilognatha as the first group of what are now the Diplopoda, and in 1840 the German naturalist Johann Friedrich von Brandt produced the first detailed classification. The name Diplopoda itself was coined in 1844 by the French zoologist Henri Marie Ducrotay de Blainville. From 1890 to 1940, millipede taxonomy was driven by relatively few researchers at any given time, with major contributions by Carl Attems, Karl Wilhelm Verhoeff and Ralph Vary Chamberlin, who each described over 1,000 species, as well as Orator F. Cook, Filippo Silvestri, R. I. Pocock, and Henry W. Brölemann. This was a period when the science of diplopodology flourished: rates of species descriptions were on average the highest in history, sometimes exceeding 300 per year.
In 1971, the Dutch biologist C. A. W. Jeekel published a comprehensive listing of all known millipede genera and families described between 1758 and 1957 in his Nomenclator Generum et Familiarum Diplopodorum, a work credited as launching the "modern era" of millipede taxonomy. In 1980, the American biologist Richard L. Hoffman published a classification of millipedes which recognized the Penicillata, Pentazonia, and Helminthomorpha, and the first phylogenetic analysis of millipede orders using modern cladistic methods was published in 1984 by Henrik Enghoff of Denmark. A 2003 classification by the American myriapodologist Rowland Shelley is similar to the one originally proposed by Verhoeff, and remains the currently accepted classification scheme (shown below), despite more recent molecular studies proposing conflicting relationships. A 2011 summary of millipede family diversity by William A. Shear placed the order Siphoniulida within the larger group Nematophora.
Fossil record
In addition to the 16 living orders, there are 9 extinct orders and one superfamily known only from fossils. The relationship of these to living groups and to each other is controversial. The extinct Arthropleuridea was long considered a distinct myriapod class, although work in the early 21st century established the group as a subclass of millipedes. Several living orders also appear in the fossil record. Below are two proposed arrangements of fossil millipede groups. Extinct groups are indicated with a dagger (†). The extinct order Zosterogrammida, a chilognath of uncertain position, is not shown.
Relation to other myriapods
Although the relationships of millipede orders are still the subject of debate, the class Diplopoda as a whole is considered a monophyletic group of arthropods: all millipedes are more closely related to each other than to any other arthropods. Diplopoda is a class within the arthropod subphylum Myriapoda, the myriapods, which includes centipedes (class Chilopoda) as well as the lesser-known pauropods (class Pauropoda) and symphylans (class Symphyla). Within myriapods, the closest relatives or sister group of millipedes has long been considered the pauropods, which also have a collum and diplosegments.
Distinction from centipedes
The differences between millipedes and centipedes are a common question from the general public. Both groups of myriapods share similarities, such as long, multi-segmented bodies, many legs, a single pair of antennae, and the presence of postantennal organs, but have many differences and distinct evolutionary histories, as the most recent common ancestor of centipedes and millipedes lived around 450 to 475 million years ago in the Silurian. The head alone exemplifies the differences; millipedes have short, geniculate (elbowed) antennae for probing the substrate, a pair of robust mandibles and a single pair of maxillae fused into a lip; centipedes have long, threadlike antennae, a pair of small mandibles, two pairs of maxillae and a pair of large poison claws.
Characteristics
Millipedes come in a variety of body shapes and sizes, ranging from to around in length, and can have as few as eleven to over three hundred segments. They are generally black or brown in colour, although there are a few brightly coloured species, and some have aposematic colouring to warn that they are toxic. Species of Motyxia produce cyanide as a chemical defence and are bioluminescent.
Body styles vary greatly between major millipede groups. In the basal subclass Penicillata, consisting of the tiny bristle millipedes, the exoskeleton is soft and uncalcified, and is covered in prominent setae or bristles. All other millipedes, belonging to the subclass Chilognatha, have a hardened exoskeleton. The chilognaths are in turn divided into two infraclasses: the Pentazonia, containing relatively short-bodied groups such as pill millipedes, and the Helminthomorpha ("worm-like" millipedes), which contains the vast majority of species, with long, many-segmented bodies.
They have also lost the gene that codes for the JHAMTl enzyme, which is responsible for catalysing the last step of the production of a juvenile hormone that regulates the development and reproduction in other arthropods like crustaceans, centipedes and insects.
Head
The head of a millipede is typically rounded above and flattened below and bears a pair of large mandibles in front of a plate-like structure called a gnathochilarium ("jaw lip"). The head contains a single pair of antennae with seven or eight segments and a group of sensory cones at the tip. Many orders also possess a pair of sensory organs known as the Tömösváry organs, shaped as small oval rings posterior and lateral to the base of the antennae. Their function is unknown, but they also occur in some centipedes, and are possibly used to measure humidity or light levels in the surrounding environment.
Millipede eyes consist of several simple flat-lensed ocelli arranged in a group or patch on each side of the head. These patches are also called ocular fields or ocellaria. Many species of millipedes, including the entire orders Polydesmida, Siphoniulida, Glomeridesmida, Siphonophorida and Platydesmida, and cave-dwelling millipedes such as Causeyella and Trichopetalum, had ancestors that could see but have subsequently lost their eyes and are blind.
Body
Millipede bodies may be flattened or cylindrical, and are composed of numerous metameric segments, each with an exoskeleton consisting of four chitinous plates: a single plate above (the tergite), one at each side (pleurites), and a plate on the underside (sternite) where the legs attach. In many millipedes, such as Merocheta and Juliformia, these plates are fused to varying degrees, sometimes forming a single cylindrical ring. The plates are typically hard, impregnated with calcium salts. Because they can't close their permanently open spiracles and most species lack a waxy cuticle, millipedes are susceptible to water loss and with a few exceptions must spend most of their time in moist or humid environments.
The first segment behind the head is legless and known as a collum (from the Latin for neck or collar). The second, third, and fourth body segments bear a single pair of legs each and are known as "haplosegments" (the three haplosegments are sometimes referred to as a "thorax"). The remaining segments, from the fifth to the posterior, are properly known as diplosegments or double segments, formed by the fusion of two embryonic segments. Each diplosegment bears two pairs of legs, rather than just one as in centipedes. In some millipedes, the last few segments may be legless. The terms "segment" or "body ring" are often used interchangeably to refer to both haplo- and diplosegments. The final segment is known as the telson and consists of a legless preanal ring, a pair of anal valves (closeable plates around the anus), and a small scale below the anus.
Millipedes in several orders have keel-like extensions of the body-wall known as paranota, which can vary widely in shape, size, and texture; modifications include lobes, papillae, ridges, crests, spines and notches. Paranota may allow millipedes to wedge more securely into crevices, protect the legs, or make the millipede more difficult for predators to swallow.
The legs are composed of seven segments, and attach on the underside of the body. The legs of an individual are generally rather similar to each other, although often longer in males than females, and males of some species may have a reduced or enlarged first pair of legs. The most conspicuous leg modifications are involved in reproduction, discussed below. Despite the common name, no millipede was known to have 1,000 legs until 2021: common species have between 34 and 400 legs, and the record is held by Eumillipes persephone, with individuals possessing up to 1,306 legs – more than any other creature on Earth.
Internal organs
Millipedes breathe through two pairs of spiracles located ventrally on each segment near the base of the legs. Each opens into an internal pouch, and connects to a system of tracheae. The heart runs the entire length of the body, with an aorta stretching into the head. The excretory organs are two pairs of malpighian tubules, located near the mid-part of the gut. The digestive tract is a simple tube with two pairs of salivary glands to help digest the food.
Reproduction and growth
Millipedes show a diversity of mating styles and structures. In the basal order Polyxenida (bristle millipedes), mating is indirect: males deposit spermatophores onto webs they secrete with special glands, and the spermatophores are subsequently picked up by females. In all other millipede groups, males possess one or two pairs of modified legs called gonopods which are used to transfer sperm to the female during copulation. The location of the gonopods differs between groups: in males of the Pentazonia they are located at the rear of the body and known as telopods and may also function in grasping females, while in the Helminthomorpha – the vast majority of species – they are located on the seventh body segment. A few species are parthenogenetic, having few, if any, males.
Gonopods occur in a diversity of shapes and sizes, and in the range from closely resembling walking legs to complex structures quite unlike legs at all. In some groups, the gonopods are kept retracted within the body; in others they project forward parallel to the body. Gonopod morphology is the predominant means of determining species among millipedes: the structures may differ greatly between closely related species but very little within a species. The gonopods develop gradually from walking legs through successive moults until reproductive maturity.
The genital openings (gonopores) of both sexes are located on the underside of the third body segment (near the second pair of legs) and may be accompanied in the male by one or two penes which deposit the sperm packets onto the gonopods. In the female, the genital pores open into paired small sacs called cyphopods or vulvae, which are covered by small hood-like lids, and are used to store the sperm after copulation. The cyphopod morphology can also be used to identify species. Millipede sperm lack flagella, a unique trait among myriapods.
In all except the bristle millipedes, copulation occurs with the two individuals facing one another. Copulation may be preceded by male behaviours such as tapping with antennae, running along the back of the female, offering edible glandular secretions, or in the case of some pill-millipedes, stridulation or "chirping". During copulation in most millipedes, the male positions his seventh segment in front of the female's third segment, and may insert his gonopods to extrude the vulvae before bending his body to deposit sperm onto his gonopods and reinserting the "charged" gonopods into the female.
Females lay from ten to three hundred eggs at a time, depending on species, fertilising them with the stored sperm as they do so. Many species deposit the eggs on moist soil or organic detritus, but some construct nests lined with dried faeces, and may protect the eggs within silk cocoons. In most species, the female abandons the eggs after they are laid, but some species in the orders Platydesmida and Stemmiulida provide parental care for eggs and young.
The young hatch after a few weeks, and typically have only three pairs of legs, followed by up to four legless segments. As they grow, they continually moult, adding further segments and legs as they do so, a mode of development known as anamorphosis. Some species moult within specially prepared chambers of soil or silk, and may also shelter in these during wet weather, and most species eat the discarded exoskeleton after moulting. The adult stage, when individuals become reproductively mature, is generally reached in the final moult stage, which varies between species and orders, although some species continue to moult after adulthood. Furthermore, some species alternate between reproductive and non-reproductive stages after maturity, a phenomenon known as periodomorphosis, in which the reproductive structures regress during non-reproductive stages. Millipedes may live from one to ten years, depending on species.
Ecology
Habitat and distribution
Millipedes occur on all continents except Antarctica, and occupy almost all terrestrial habitats, ranging as far north as the Arctic Circle in Iceland, Norway, and Central Russia, and as far south as Santa Cruz Province, Argentina. Typically forest floor dwellers, they live in leaf litter, dead wood, or soil, with a preference for humid conditions. In temperate zones, millipedes are most abundant in moist deciduous forests, and may reach densities of over 1,000 individuals per square metre. Other habitats include coniferous forests, caves, and alpine ecosystems. Deserticolous millipedes, species evolved to live in the desert, like Orthoporus ornatus, may show adaptations like a waxy epicuticle and the ability of water uptake from unsaturated air. Some species can survive freshwater floods and live submerged underwater for up to 11 months. A few species occur near the seashore and can survive in somewhat salty conditions.
Burrowing
The diplosegments of millipedes have evolved in conjunction with their burrowing habits, and nearly all millipedes adopt a mainly subterranean lifestyle. They use three main methods of burrowing; bulldozing, wedging and boring. Members of the orders Julida, Spirobolida and Spirostreptida, lower their heads and barge their way into the substrate, the collum leading the way. Flat-backed millipedes in the order Polydesmida tend to insert their front end, like a wedge, into a horizontal crevice, and then widen the crack by pushing upwards with their legs, the paranota in this instance constituting the main lifting surface. Boring is used by members of the order Polyzoniida. These have smaller segments at the front and increasingly large ones further back; they propel themselves forward into a crack with their legs, the wedge-shaped body widening the gap as they go. Some millipedes have adopted an above-ground lifestyle and lost the burrowing habit. This may be because they are too small to have enough leverage to burrow, or because they are too large to make the effort worthwhile, or in some cases because they move relatively fast (for a millipede) and are active predators.
Diet
Most millipedes are detritivores and feed on decomposing vegetation, feces, or organic matter mixed with soil. They often play important roles in the breakdown and decomposition of plant litter: estimates of consumption rates for individual species range from 1 to 11 percent of all leaf litter, depending on species and region, and collectively millipedes may consume nearly all the leaf litter in a region. The leaf litter is fragmented in the millipede gut and excreted as pellets of leaf fragments, algae, fungi, and bacteria, which facilitates decomposition by the microorganisms. Where earthworm populations are low in tropical forests, millipedes play an important role in facilitating microbial decomposition of the leaf litter. Some millipedes are herbivorous, feeding on living plants, and some species can become serious pests of crops. Millipedes in the order Polyxenida graze algae from bark, and Platydesmida feed on fungi. A few species are omnivorous or in Callipodida and Chordeumatida occasionally carnivorous, feeding on insects, centipedes, earthworms, or snails. Some species have piercing mouth parts that allow them to suck up plant juices. Cave dwelling species in Julidae, Blaniulidae, and Polydesmidae have specialized mouthparts and appears to be filter feeders, filtering small particles from running water inside caves.
Predators and parasites
Millipedes are preyed on by a wide range of animals, including various reptiles, amphibians, birds, mammals, and insects. Mammalian predators such as coatis and meerkats roll captured millipedes on the ground to deplete and rub off their defensive secretions before consuming their prey, and certain poison dart frogs are believed to incorporate the toxic compounds of millipedes into their own defences. Several invertebrates have specialised behaviours or structures to feed on millipedes, including larval glowworm beetles, Probolomyrmex ants, chlamydephorid slugs, and predaceous dung beetles of the genera Sceliages and Deltochilum. A large subfamily of assassin bugs, the Ectrichodiinae with over 600 species, has specialised in preying upon millipedes. Parasites of millipedes include nematodes, phaeomyiid flies, and acanthocephalans. Nearly 30 fungal species of the order Laboulbeniales have been found growing externally on millipedes, but some species may be commensal rather than parasitic.
Defence mechanisms
Due to their lack of speed and their inability to bite or sting, millipedes' primary defence mechanism is to curl into a tight coil – protecting their delicate legs inside an armoured exoskeleton.
Many species also emit various foul-smelling liquid secretions through microscopic holes called ozopores (the openings of "odoriferous" or "repugnatorial glands"), along the sides of their bodies as a secondary defence. Among the many irritant and toxic chemicals found in these secretions are alkaloids, benzoquinones, phenols, terpenoids, and hydrogen cyanide. Some of these substances are caustic and can burn the exoskeleton of ants and other insect predators, and the skin and eyes of larger predators. Primates such as capuchin monkeys and lemurs have been observed intentionally irritating millipedes in order to rub the chemicals on themselves to repel mosquitoes. Some of these defensive compounds also show antifungal activity.
The bristly millipedes (order Polyxenida) lack both an armoured exoskeleton and odiferous glands, and instead are covered in numerous bristles that in at least one species, Polyxenus fasciculatus, detach and entangle ants.
Other inter-species interactions
Some millipedes form mutualistic relationships with organisms of other species, in which both species benefit from the interaction, or commensal relationships, in which only one species benefits while the other is unaffected. Several species form close relationships with ants, a relationship known as myrmecophily, especially within the family Pyrgodesmidae (Polydesmida), which contains "obligate myrmecophiles", species which have only been found in ant colonies. More species are "facultative myrmecophiles", non-exclusively associated with ants, including many species of Polyxenida that have been found in ant nests around the world.
Many millipede species have commensal relationships with mites of the orders Mesostigmata and Astigmata. Many of these mites are believed to be phoretic rather than parasitic, which means that they use the millipede host as a means of dispersal.
A novel interaction between millipedes and mosses was described in 2011, in which individuals of the newly discovered Psammodesmus bryophorus was found to have up to ten species living on its dorsal surface, in what may provide camouflage for the millipede and increased dispersal for the mosses.
Interactions with humans
Millipedes generally have little impact on human economic or social well-being, especially in comparison with insects, although locally they can be a nuisance or agricultural pest. Millipedes do not bite, and their defensive secretions are mostly harmless to humans — usually causing only minor discolouration on the skin — but the secretions of some tropical species may cause pain, itching, local erythema, edema, blisters, eczema, and occasionally cracked skin. Eye exposures to these secretions causes general irritation and potentially more severe effects such as conjunctivitis and keratitis. This is called millipede burn. First aid consists of flushing the area thoroughly with water; further treatment is aimed at relieving the local effects.
Some millipedes are considered household pests, including Xenobolus carnifex which can infest thatched roofs in India, and Ommatoiulus moreleti, which periodically invades homes in Australia. Other species exhibit periodical swarming behaviour, which can result in home invasions, crop damage, and train delays when the tracks become slippery with the crushed remains of hundreds of millipedes. Some millipedes can cause significant damage to crops: the spotted snake millipede (Blaniulus guttulatus) is a pest of sugar beets and other root crops, and as a result is one of the few millipedes with a common name.
Some of the larger millipedes in the orders Spirobolida, Spirostreptida, and Sphaerotheriida are popular as pets. Some species commonly sold or kept include species of Archispirostreptus, Aphistogoniulus, Narceus, and Orthoporus.
Millipedes appear in folklore and traditional medicine around the world. Some cultures associate millipede activity with coming rains. In Zambia, smashed millipede pulp is used to treat wounds, and the Bafia people of Cameroon use millipede juice to treat earache. In certain Himalayan Bhotiya tribes, dry millipede smoke is used to treat haemorrhoids. Native people in Malaysia use millipede secretions in poison-tipped arrows. The secretions of Spirobolus bungii have been observed to inhibit division of human cancer cells. The only recorded usage of millipedes as food by humans comes from the Bobo people of Burkina Faso in West Africa, who consume boiled, dried millipedes belonging to the families Gomphodesmidae and Spirostreptidae to which they add tomato sauce.
Millipedes have also inspired and played roles in scientific research. In 1963, a walking vehicle with 36 legs was designed, said to have been inspired by a study of millipede locomotion.
Experimental robots have had the same inspiration, in particular when heavy loads are needed to be carried in tight areas involving turns and curves. In biology, some authors have advocated millipedes as model organisms for the study of arthropod physiology and the developmental processes controlling the number and shape of body segments.
Similar to vermicompost, millipedes can be used to convert plant matter into compost in what has been named millicomposting, which improves the quality of the compost.
| Biology and health sciences | Myriapoda | null |
100563 | https://en.wikipedia.org/wiki/System%20on%20a%20chip | System on a chip | A system on a chip or system-on-chip (SoC ; pl. SoCs ) is an integrated circuit that integrates most or all components of a computer or electronic system. These components usually include an on-chip central processing unit (CPU), memory interfaces, input/output devices and interfaces, and secondary storage interfaces, often alongside other components such as radio modems and a graphics processing unit (GPU) – all on a single substrate or microchip. SoCs may contain digital and also analog, mixed-signal and often radio frequency signal processing functions (otherwise it may be considered on a discrete application processor).
High-performance SoCs are often paired with dedicated and physically separate memory and secondary storage (such as LPDDR and eUFS or eMMC, respectively) chips that may be layered on top of the SoC in what is known as a package on package (PoP) configuration, or be placed close to the SoC. Additionally, SoCs may use separate wireless modems (especially WWAN modems).
An SoC integrates a microcontroller, microprocessor or perhaps several processor cores with peripherals like a GPU, Wi-Fi and cellular network radio modems or one or more coprocessors. Similar to how a microcontroller integrates a microprocessor with peripheral circuits and memory, an SoC can be seen as integrating a microcontroller with even more advanced peripherals.
Compared to a multi-chip architecture, an SoC with equivalent functionality will have reduced power consumption as well as a smaller semiconductor die area. This comes at the cost of reduced replaceability of components. By definition, SoC designs are fully or nearly fully integrated across different component modules. For these reasons, there has been a general trend towards tighter integration of components in the computer hardware industry, in part due to the influence of SoCs and lessons learned from the mobile and embedded computing markets.
SoCs are very common in the mobile computing (as in smart devices such as smartphones and tablet computers) and edge computing markets.
Types
In general, there are three distinguishable types of SoCs:
SoCs built around a microcontroller,
SoCs built around a microprocessor, often found in mobile phones;
Specialized application-specific integrated circuit SoCs designed for specific applications that do not fit into the above two categories.
Applications
SoCs can be applied to any computing task. However, they are typically used in mobile computing such as tablets, smartphones, smartwatches, and netbooks as well as embedded systems and in applications where previously microcontrollers would be used.
Embedded systems
Where previously only microcontrollers could be used, SoCs are rising to prominence in the embedded systems market. Tighter system integration offers better reliability and mean time between failure, and SoCs offer more advanced functionality and computing power than microcontrollers. Applications include AI acceleration, embedded machine vision, data collection, telemetry, vector processing and ambient intelligence. Often embedded SoCs target the internet of things, multimedia, networking, telecommunications and edge computing markets. Some examples of SoCs for embedded applications include:
AMD
Zynq 7000 SoC
Zynq UltraScale+ MPSoC
Zynq UltraScale+ RFSoC
Versal Adaptive SoC
Mobile computing
Mobile computing based SoCs always bundle processors, memories, on-chip caches, wireless networking capabilities and often digital camera hardware and firmware. With increasing memory sizes, high end SoCs will often have no memory and flash storage and instead, the memory and flash memory will be placed right next to, or above (package on package), the SoC. Some examples of mobile computing SoCs include:
Samsung Electronics: list, typically based on ARM
Exynos, used mainly by Samsung's Galaxy series of smartphones
Qualcomm:
Snapdragon (list), used in many smartphones. In 2018, Snapdragon SoCs were being used as the backbone of laptop computers running Windows 10, marketed as "Always Connected PCs".
MediaTek, typically based on ARM
Dimensity & Kompanio Series. Standalone application & tablet processors that power devices such as Amazon Echo Show
Personal computers
In 1992, Acorn Computers produced the A3010, A3020 and A4000 range of personal computers with the ARM250 SoC. It combined the original Acorn ARM2 processor with a memory controller (MEMC), video controller (VIDC), and I/O controller (IOC). In previous Acorn ARM-powered computers, these were four discrete chips. The ARM7500 chip was their second-generation SoC, based on the ARM700, VIDC20 and IOMD controllers, and was widely licensed in embedded devices such as set-top-boxes, as well as later Acorn personal computers.
Tablet and laptop manufacturers have learned lessons from embedded systems and smartphone markets about reduced power consumption, better performance and reliability from tighter integration of hardware and firmware modules, and LTE and other wireless network communications integrated on chip (integrated network interface controllers).
On modern laptops and mini PCs, the low-power variants of AMD Ryzen and Intel Core processors, are use SoC design integrating CPU, IGPU, chipset and other processors in a single package. However, such x86 processors still require external memory and storage chips.
Structure
An SoC consists of hardware functional units, including microprocessors that run software code, as well as a communications subsystem to connect, control, direct and interface between these functional modules.
Functional components
Processor cores
An SoC must have at least one processor core, but typically an SoC has more than one core. Processor cores can be a microcontroller, microprocessor (μP), digital signal processor (DSP) or application-specific instruction set processor (ASIP) core. ASIPs have instruction sets that are customized for an application domain and designed to be more efficient than general-purpose instructions for a specific type of workload. Multiprocessor SoCs have more than one processor core by definition. The ARM architecture is a common choice for SoC processor cores because some ARM-architecture cores are soft processors specified as IP cores.
Memory
SoCs must have semiconductor memory blocks to perform their computation, as do microcontrollers and other embedded systems. Depending on the application, SoC memory may form a memory hierarchy and cache hierarchy. In the mobile computing market, this is common, but in many low-power embedded microcontrollers, this is not necessary. Memory technologies for SoCs include read-only memory (ROM), random-access memory (RAM), Electrically Erasable Programmable ROM (EEPROM) and flash memory. As in other computer systems, RAM can be subdivided into relatively faster but more expensive static RAM (SRAM) and the slower but cheaper dynamic RAM (DRAM). When an SoC has a cache hierarchy, SRAM will usually be used to implement processor registers and cores' built-in caches whereas DRAM will be used for main memory. "Main memory" may be specific to a single processor (which can be multi-core) when the SoC has multiple processors, in this case it is distributed memory and must be sent via on-chip to be accessed by a different processor. For further discussion of multi-processing memory issues, see cache coherence and memory latency.
Interfaces
SoCs include external interfaces, typically for communication protocols. These are often based upon industry standards such as USB, Ethernet, USART, SPI, HDMI, I²C, CSI, etc. These interfaces will differ according to the intended application. Wireless networking protocols such as Wi-Fi, Bluetooth, 6LoWPAN and near-field communication may also be supported.
When needed, SoCs include analog interfaces including analog-to-digital and digital-to-analog converters, often for signal processing. These may be able to interface with different types of sensors or actuators, including smart transducers. They may interface with application-specific modules or shields. Or they may be internal to the SoC, such as if an analog sensor is built in to the SoC and its readings must be converted to digital signals for mathematical processing.
Digital signal processors
Digital signal processor (DSP) cores are often included on SoCs. They perform signal processing operations in SoCs for sensors, actuators, data collection, data analysis and multimedia processing. DSP cores typically feature very long instruction word (VLIW) and single instruction, multiple data (SIMD) instruction set architectures, and are therefore highly amenable to exploiting instruction-level parallelism through parallel processing and superscalar execution. SP cores most often feature application-specific instructions, and as such are typically application-specific instruction set processors (ASIP). Such application-specific instructions correspond to dedicated hardware functional units that compute those instructions.
Typical DSP instructions include multiply-accumulate, Fast Fourier transform, fused multiply-add, and convolutions.
Other
As with other computer systems, SoCs require timing sources to generate clock signals, control execution of SoC functions and provide time context to signal processing applications of the SoC, if needed. Popular time sources are crystal oscillators and phase-locked loops.
SoC peripherals including counter-timers, real-time timers and power-on reset generators. SoCs also include voltage regulators and power management circuits.
Intermodule communication
SoCs comprise many execution units. These units must often send data and instructions back and forth. Because of this, all but the most trivial SoCs require communications subsystems. Originally, as with other microcomputer technologies, data bus architectures were used, but recently designs based on sparse intercommunication networks known as networks-on-chip (NoC) have risen to prominence and are forecast to overtake bus architectures for SoC design in the near future.
Bus-based communication
Historically, a shared global computer bus typically connected the different components, also called "blocks" of the SoC. A very common bus for SoC communications is ARM's royalty-free Advanced Microcontroller Bus Architecture (AMBA) standard.
Direct memory access controllers route data directly between external interfaces and SoC memory, bypassing the CPU or control unit, thereby increasing the data throughput of the SoC. This is similar to some device drivers of peripherals on component-based multi-chip module PC architectures.
Wire delay is not scalable due to continued miniaturization, system performance does not scale with the number of cores attached, the SoC's operating frequency must decrease with each additional core attached for power to be sustainable, and long wires consume large amounts of electrical power. These challenges are prohibitive to supporting manycore systems on chip.
Network on a chip
In the late 2010s, a trend of SoCs implementing communications subsystems in terms of a network-like topology instead of bus-based protocols has emerged. A trend towards more processor cores on SoCs has caused on-chip communication efficiency to become one of the key factors in determining the overall system performance and cost. This has led to the emergence of interconnection networks with router-based packet switching known as "networks on chip" (NoCs) to overcome the bottlenecks of bus-based networks.
Networks-on-chip have advantages including destination- and application-specific routing, greater power efficiency and reduced possibility of bus contention. Network-on-chip architectures take inspiration from communication protocols like TCP and the Internet protocol suite for on-chip communication, although they typically have fewer network layers. Optimal network-on-chip network architectures are an ongoing area of much research interest. NoC architectures range from traditional distributed computing network topologies such as torus, hypercube, meshes and tree networks to genetic algorithm scheduling to randomized algorithms such as random walks with branching and randomized time to live (TTL).
Many SoC researchers consider NoC architectures to be the future of SoC design because they have been shown to efficiently meet power and throughput needs of SoC designs. Current NoC architectures are two-dimensional. 2D IC design has limited floorplanning choices as the number of cores in SoCs increase, so as three-dimensional integrated circuits (3DICs) emerge, SoC designers are looking towards building three-dimensional on-chip networks known as 3DNoCs.
Design flow
A system on a chip consists of both the hardware, described in , and the software controlling the microcontroller, microprocessor or digital signal processor cores, peripherals and interfaces. The design flow for an SoC aims to develop this hardware and software at the same time, also known as architectural co-design. The design flow must also take into account optimizations () and constraints.
Most SoCs are developed from pre-qualified hardware component IP core specifications for the hardware elements and execution units, collectively "blocks", described above, together with software device drivers that may control their operation. Of particular importance are the protocol stacks that drive industry-standard interfaces like USB. The hardware blocks are put together using computer-aided design tools, specifically electronic design automation tools; the software modules are integrated using a software integrated development environment.
SoCs components are also often designed in high-level programming languages such as C++, MATLAB or SystemC and converted to RTL designs through high-level synthesis (HLS) tools such as C to HDL or flow to HDL. HLS products called "algorithmic synthesis" allow designers to use C++ to model and synthesize system, circuit, software and verification levels all in one high level language commonly known to computer engineers in a manner independent of time scales, which are typically specified in HDL. Other components can remain software and be compiled and embedded onto soft-core processors included in the SoC as modules in HDL as IP cores.
Once the architecture of the SoC has been defined, any new hardware elements are written in an abstract hardware description language termed register transfer level (RTL) which defines the circuit behavior, or synthesized into RTL from a high level language through high-level synthesis. These elements are connected together in a hardware description language to create the full SoC design. The logic specified to connect these components and convert between possibly different interfaces provided by different vendors is called glue logic.
Design verification
Chips are verified for validation correctness before being sent to a semiconductor foundry. This process is called functional verification and it accounts for a significant portion of the time and energy expended in the chip design life cycle, often quoted as 70%. With the growing complexity of chips, hardware verification languages like SystemVerilog, SystemC, e, and OpenVera are being used. Bugs found in the verification stage are reported to the designer.
Traditionally, engineers have employed simulation acceleration, emulation or prototyping on reprogrammable hardware to verify and debug hardware and software for SoC designs prior to the finalization of the design, known as tape-out. Field-programmable gate arrays (FPGAs) are favored for prototyping SoCs because FPGA prototypes are reprogrammable, allow debugging and are more flexible than application-specific integrated circuits (ASICs).
With high capacity and fast compilation time, simulation acceleration and emulation are powerful technologies that provide wide visibility into systems. Both technologies, however, operate slowly, on the order of MHz, which may be significantly slower – up to 100 times slower – than the SoC's operating frequency. Acceleration and emulation boxes are also very large and expensive at over US$1 million.
FPGA prototypes, in contrast, use FPGAs directly to enable engineers to validate and test at, or close to, a system's full operating frequency with real-world stimuli. Tools such as Certus are used to insert probes in the FPGA RTL that make signals available for observation. This is used to debug hardware, firmware and software interactions across multiple FPGAs with capabilities similar to a logic analyzer.
In parallel, the hardware elements are grouped and passed through a process of logic synthesis, during which performance constraints, such as operational frequency and expected signal delays, are applied. This generates an output known as a netlist describing the design as a physical circuit and its interconnections. These netlists are combined with the glue logic connecting the components to produce the schematic description of the SoC as a circuit which can be printed onto a chip. This process is known as place and route and precedes tape-out in the event that the SoCs are produced as application-specific integrated circuits (ASIC).
Optimization goals
SoCs must optimize power use, area on die, communication, positioning for locality between modular units and other factors. Optimization is necessarily a design goal of SoCs. If optimization was not necessary, the engineers would use a multi-chip module architecture without accounting for the area use, power consumption or performance of the system to the same extent.
Common optimization targets for SoC designs follow, with explanations of each. In general, optimizing any of these quantities may be a hard combinatorial optimization problem, and can indeed be NP-hard fairly easily. Therefore, sophisticated optimization algorithms are often required and it may be practical to use approximation algorithms or heuristics in some cases. Additionally, most SoC designs contain multiple variables to optimize simultaneously, so Pareto efficient solutions are sought after in SoC design. Oftentimes the goals of optimizing some of these quantities are directly at odds, further adding complexity to design optimization of SoCs and introducing trade-offs in system design.
For broader coverage of trade-offs and requirements analysis, see requirements engineering.
Targets
Power consumption
SoCs are optimized to minimize the electrical power used to perform the SoC's functions. Most SoCs must use low power. SoC systems often require long battery life (such as smartphones), can potentially spend months or years without a power source while needing to maintain autonomous function, and often are limited in power use by a high number of embedded SoCs being networked together in an area. Additionally, energy costs can be high and conserving energy will reduce the total cost of ownership of the SoC. Finally, waste heat from high energy consumption can damage other circuit components if too much heat is dissipated, giving another pragmatic reason to conserve energy. The amount of energy used in a circuit is the integral of power consumed with respect to time, and the average rate of power consumption is the product of current by voltage. Equivalently, by Ohm's law, power is current squared times resistance or voltage squared divided by resistance:
SoCs are frequently embedded in portable devices such as smartphones, GPS navigation devices, digital watches (including smartwatches) and netbooks. Customers want long battery lives for mobile computing devices, another reason that power consumption must be minimized in SoCs. Multimedia applications are often executed on these devices, including video games, video streaming, image processing; all of which have grown in computational complexity in recent years with user demands and expectations for higher-quality multimedia. Computation is more demanding as expectations move towards 3D video at high resolution with multiple standards, so SoCs performing multimedia tasks must be computationally capable platform while being low power to run off a standard mobile battery.
Performance per watt
SoCs are optimized to maximize power efficiency in performance per watt: maximize the performance of the SoC given a budget of power usage. Many applications such as edge computing, distributed processing and ambient intelligence require a certain level of computational performance, but power is limited in most SoC environments.
Waste heat
SoC designs are optimized to minimize waste heat output on the chip. As with other integrated circuits, heat generated due to high power density are the bottleneck to further miniaturization of components. The power densities of high speed integrated circuits, particularly microprocessors and including SoCs, have become highly uneven. Too much waste heat can damage circuits and erode reliability of the circuit over time. High temperatures and thermal stress negatively impact reliability, stress migration, decreased mean time between failures, electromigration, wire bonding, metastability and other performance degradation of the SoC over time.
In particular, most SoCs are in a small physical area or volume and therefore the effects of waste heat are compounded because there is little room for it to diffuse out of the system. Because of high transistor counts on modern devices, oftentimes a layout of sufficient throughput and high transistor density is physically realizable from fabrication processes but would result in unacceptably high amounts of heat in the circuit's volume.
These thermal effects force SoC and other chip designers to apply conservative design margins, creating less performant devices to mitigate the risk of catastrophic failure. Due to increased transistor densities as length scales get smaller, each process generation produces more heat output than the last. Compounding this problem, SoC architectures are usually heterogeneous, creating spatially inhomogeneous heat fluxes, which cannot be effectively mitigated by uniform passive cooling.
Throughput
SoCs are optimized to maximize computational and communications throughput.
Latency
SoCs are optimized to minimize latency for some or all of their functions. This can be accomplished by laying out elements with proper proximity and locality to each-other to minimize the interconnection delays and maximize the speed at which data is communicated between modules, functional units and memories. In general, optimizing to minimize latency is an NP-complete problem equivalent to the Boolean satisfiability problem.
For tasks running on processor cores, latency and throughput can be improved with task scheduling. Some tasks run in application-specific hardware units, however, and even task scheduling may not be sufficient to optimize all software-based tasks to meet timing and throughput constraints.
Methodologies
Systems on chip are modeled with standard hardware verification and validation techniques, but additional techniques are used to model and optimize SoC design alternatives to make the system optimal with respect to multiple-criteria decision analysis on the above optimization targets.
Task scheduling
Task scheduling is an important activity in any computer system with multiple processes or threads sharing a single processor core. It is important to reduce and increase for embedded software running on an SoC's . Not every important computing activity in a SoC is performed in software running on on-chip processors, but scheduling can drastically improve performance of software-based tasks and other tasks involving shared resources.
Software running on SoCs often schedules tasks according to network scheduling and randomized scheduling algorithms.
Pipelining
Hardware and software tasks are often pipelined in processor design. Pipelining is an important principle for speedup in computer architecture. They are frequently used in GPUs (graphics pipeline) and RISC processors (evolutions of the classic RISC pipeline), but are also applied to application-specific tasks such as digital signal processing and multimedia manipulations in the context of SoCs.
Probabilistic modeling
SoCs are often analyzed though probabilistic models, queueing networks, and Markov chains. For instance, Little's law allows SoC states and NoC buffers to be modeled as arrival processes and analyzed through Poisson random variables and Poisson processes.
Markov chains
SoCs are often modeled with Markov chains, both discrete time and continuous time variants. Markov chain modeling allows asymptotic analysis of the SoC's steady state distribution of power, heat, latency and other factors to allow design decisions to be optimized for the common case.
Fabrication
SoC chips are typically fabricated using metal–oxide–semiconductor (MOS) technology. The netlists described above are used as the basis for the physical design (place and route) flow to convert the designers' intent into the design of the SoC. Throughout this conversion process, the design is analyzed with static timing modeling, simulation and other tools to ensure that it meets the specified operational parameters such as frequency, power consumption and dissipation, functional integrity (as described in the register transfer level code) and electrical integrity.
When all known bugs have been rectified and these have been re-verified and all physical design checks are done, the physical design files describing each layer of the chip are sent to the foundry's mask shop where a full set of glass lithographic masks will be etched. These are sent to a wafer fabrication plant to create the SoC dice before packaging and testing.
SoCs can be fabricated by several technologies, including:
Full custom ASIC
Standard cell ASIC
Field-programmable gate array (FPGA)
ASICs consume less power and are faster than FPGAs but cannot be reprogrammed and are expensive to manufacture. FPGA designs are more suitable for lower volume designs, but after enough units of production ASICs reduce the total cost of ownership.
SoC designs consume less power and have a lower cost and higher reliability than the multi-chip systems that they replace. With fewer packages in the system, assembly costs are reduced as well.
However, like most very-large-scale integration (VLSI) designs, the total cost is higher for one large chip than for the same functionality distributed over several smaller chips, because of lower yields and higher non-recurring engineering costs.
When it is not feasible to construct an SoC for a particular application, an alternative is a system in package (SiP) comprising a number of chips in a single package. When produced in large volumes, SoC is more cost-effective than SiP because its packaging is simpler. Another reason SiP may be preferred is waste heat may be too high in a SoC for a given purpose because functional components are too close together, and in an SiP heat will dissipate better from different functional modules since they are physically further apart.
Examples
Some examples of systems on a chip are:
Apple A series
Cell processor
Adapteva's Epiphany architecture
Xilinx Zynq UltraScale
Qualcomm Snapdragon
Benchmarks
SoC research and development often compares many options. Benchmarks, such as COSMIC, are developed to help such evaluations.
| Technology | Semiconductors | null |
100625 | https://en.wikipedia.org/wiki/Inverter%20%28logic%20gate%29 | Inverter (logic gate) | In digital logic, an inverter or NOT gate is a logic gate which implements logical negation. It outputs a bit opposite of the bit that is put into it. The bits are typically implemented as two differing voltage levels.
Description
The NOT gate outputs a zero when given a one, and a one when given a zero. Hence, it inverts its inputs. Colloquially, this inversion of bits is called "flipping" bits. As with all binary logic gates, other pairs of symbols such as true and false, or high and low may be used in lieu of one and zero.
It is equivalent to the logical negation operator (¬) in mathematical logic. Because it has only one input, it is a unary operation and has the simplest type of truth table. It is also called the complement gate because it produces the ones' complement of a binary number, swapping 0s and 1s.
The NOT gate is one of three basic logic gates from which any Boolean circuit may be built up. Together with the AND gate and the OR gate, any function in binary mathematics may be implemented. All other logic gates may be made from these three.
The terms "programmable inverter" or "controlled inverter" do not refer to this gate; instead, these terms refer to the XOR gate because it can conditionally function like a NOT gate.
Symbols
The traditional symbol for an inverter circuit is a triangle touching a small circle or "bubble". Input and output lines are attached to the symbol; the bubble is typically attached to the output line. To symbolize active-low input, sometimes the bubble is instead placed on the input line. Sometimes only the circle portion of the symbol is used, and it is attached to the input or output of another gate; the symbols for NAND and NOR are formed in this way.
A bar or overline ( ‾ ) above a variable can denote negation (or inversion or complement) performed by a NOT gate. A slash (/) before the variable is also used.
Electronic implementation
An inverter circuit outputs a voltage representing the opposite logic-level to its input. Its main function is to invert the input signal applied. If the applied input is low then the output becomes high and vice versa. Inverters can be constructed using a single NMOS transistor or a single PMOS transistor coupled with a resistor. Since this "resistive-drain" approach uses only a single type of transistor, it can be fabricated at a low cost. However, because current flows through the resistor in one of the two states, the resistive-drain configuration is disadvantaged for power consumption and processing speed. Alternatively, inverters can be constructed using two complementary transistors in a CMOS configuration. This configuration greatly reduces power consumption since one of the transistors is always off in both logic states. Processing speed can also be improved due to the relatively low resistance compared to the NMOS-only or PMOS-only type devices. Inverters can also be constructed with bipolar junction transistors (BJT) in either a resistor–transistor logic (RTL) or a transistor–transistor logic (TTL) configuration.
Digital electronics circuits operate at fixed voltage levels corresponding to a logical 0 or 1 (see binary). An inverter circuit serves as the basic logic gate to swap between those two voltage levels. Implementation determines the actual voltage, but common levels include (0, +5V) for TTL circuits.
Digital building block
The inverter is a basic building block in digital electronics. Multiplexers, decoders, state machines, and other sophisticated digital devices may use inverters.
The hex inverter is an integrated circuit that contains six (hexa-) inverters. For example, the 7404 TTL chip which has 14 pins and the 4049 CMOS chip which has 16 pins, 2 of which are used for power/referencing, and 12 of which are used by the inputs and outputs of the six inverters (the 4049 has 2 pins with no connection).
Analytical representation
is the analytical representation of NOT gate:
Alternatives
If no specific NOT gates are available, one can be made from the universal NAND or NOR gates, or an XOR gate by setting one input to high.
Performance measurement
Digital inverter quality is often measured using the voltage transfer curve (VTC), which is a plot of output vs. input voltage. From such a graph, device parameters including noise tolerance, gain, and operating logic levels can be obtained.
Ideally, the VTC appears as an inverted step function – this would indicate precise switching between on and off – but in real devices, a gradual transition region exists. The VTC indicates that for low input voltage, the circuit outputs high voltage; for high input, the output tapers off towards the low level. The slope of this transition region is a measure of quality – steep (close to vertical) slopes yield precise switching.
The tolerance to noise can be measured by comparing the minimum input to the maximum output for each region of operation (on / off).
Linear region as analog amplifier
Since the transition region is steep and approximately linear, a properly-biased CMOS inverter digital logic gate may be used as a high-gain analog linear amplifier or even combined to form an opamp. Maximum gain is achieved when the input and output operating points are the same voltage, which can be biased by connecting a resistor between the output and input.
| Technology | Digital logic | null |
101107 | https://en.wikipedia.org/wiki/Dentition | Dentition | Dentition pertains to the development of teeth and their arrangement in the mouth. In particular, it is the characteristic arrangement, kind, and number of teeth in a given species at a given age. That is, the number, type, and morpho-physiology (that is, the relationship between the shape and form of the tooth in question and its inferred function) of the teeth of an animal.
Terminology
Animals whose teeth are all of the same type, such as most non-mammalian vertebrates, are said to have homodont dentition, whereas those whose teeth differ morphologically are said to have heterodont dentition. The dentition of animals with two successions of teeth (deciduous, permanent) is referred to as diphyodont, while the dentition of animals with only one set of teeth throughout life is monophyodont. The dentition of animals in which the teeth are continuously discarded and replaced throughout life is termed polyphyodont. The dentition of animals in which the teeth are set in sockets in the jawbones is termed thecodont.
Overview
The evolutionary origin of the vertebrate dentition remains contentious. Current theories suggest either an "outside-in" or "inside-out" evolutionary origin to teeth, with the dentition arising from odontodes on the skin surface moving into the mouth, or vice versa. Despite this debate, it is accepted that vertebrate teeth are homologous to the dermal denticles found on the skin of basal Gnathostomes (i.e. Chondrichtyans). Since the origin of teeth some 450 mya, the vertebrate dentition has diversified within the reptiles, amphibians, and fish: however most of these groups continue to possess a long row of pointed or sharp-sided, undifferentiated teeth (homodont) that are completely replaceable. The mammalian pattern is significantly different. The teeth in the upper and lower jaws in mammals have evolved a close-fitting relationship such that they operate together as a unit. "They 'occlude', that is, the chewing surfaces of the teeth are so constructed that the upper and lower teeth are able to fit precisely together, cutting, crushing, grinding or tearing the food caught between."
Mammals have up to four distinct types of teeth, though not all types are present in all mammals. These are the incisor (cutting), the canine, the premolar, and the molar (grinding). The incisors occupy the front of the tooth row in both upper and lower jaws. They are normally flat, chisel-shaped teeth that meet in an edge-to-edge bite. Their function is cutting, slicing, or gnawing food into manageable pieces that fit into the mouth for further chewing. The canines are immediately behind the incisors. In many mammals, the canines are pointed, tusk-shaped teeth, projecting beyond the level of the other teeth. In carnivores, they are primarily offensive weapons for bringing down prey. In other mammals such as some primates, they are used to split open hard-surfaced food. In humans, the canine teeth are the main components in occlusal function and articulation. The mandibular teeth function against the maxillary teeth in a particular movement that is harmonious to the shape of the occluding surfaces. This creates the incising and grinding functions. The teeth must mesh together the way gears mesh in a transmission. If the interdigitation of the opposing cusps and incisal edges are not directed properly the teeth will wear abnormally (attrition), break away irregular crystalline enamel structures from the surface (abrasion), or fracture larger pieces (abfraction). This is a three-dimensional movement of the mandible in relation to the maxilla. There are three points of guidance: the two posterior points provided by the temporomandibular joints and the anterior component provided by the incisors and canines. The incisors mostly control the vertical opening of the chewing cycle when the muscles of mastication move the jaw forwards and backwards (protrusion/retrusion). The canines come into function guiding the vertical movement when the chewing is side to side (lateral). The canines alone can cause the other teeth to separate at the extreme end of the cycle (cuspid guided function) or all the posterior teeth can continue to stay in contact (group function). The entire range of this movement is the envelope of masticatory function. The initial movement inside this envelope is directed by the shape of the teeth in contact and the Glenoid Fossa/Condyle shape. The outer extremities of this envelope are limited by muscles, ligaments and the articular disc of the TMJ. Without the guidance of anterior incisors and canines, this envelope of function can be destructive to the remaining teeth resulting in periodontal trauma from occlusion seen as wear, fracture or tooth loosening and loss. The premolars and molars are at the back of the mouth. Depending on the particular mammal and its diet, these two kinds of teeth prepare pieces of food to be swallowed by grinding, shearing, or crushing. The specialised teeth—incisors, canines, premolars, and molars—are found in the same order in every mammal. In many mammals, the infants have a set of teeth that fall out and are replaced by adult teeth. These are called deciduous teeth, primary teeth, baby teeth or milk teeth. Animals that have two sets of teeth, one followed by the other, are said to be diphyodont. Normally the dental formula for milk teeth is the same as for adult teeth except that the molars are missing.
Dental formula
Because every mammal's teeth are specialised for different functions, many mammal groups have lost the teeth that are not needed in their adaptation. Tooth form has also undergone evolutionary modification as a result of natural selection for specialised feeding or other adaptations. Over time, different mammal groups have evolved distinct dental features, both in the number and type of teeth and in the shape and size of the chewing surface.
The number of teeth of each type is written as a dental formula for one side of the mouth, or quadrant, with the upper and lower teeth shown on separate rows. The number of teeth in a mouth is twice that listed, as there are two sides. In each set, incisors (I) are indicated first, canines (C) second, premolars (P) third, and finally molars (M), giving I:C:P:M. So for example, the formula 2.1.2.3 for upper teeth indicates 2 incisors, 1 canine, 2 premolars, and 3 molars on one side of the upper mouth. The deciduous dental formula is notated in lowercase lettering preceded by the letter d: for example: di:dc:dp.
An animal's dentition for either deciduous or permanent teeth can thus be expressed as a dental formula, written in the form of a fraction, which can be written as , or I.C.P.M / I.C.P.M. For example, the following formulae show the deciduous and usual permanent dentition of all catarrhine primates, including humans:
Deciduous: This can also be written as . Superscript and subscript denote upper and lower jaw, i.e. do not indicate mathematical operations; the numbers are the count of the teeth of each type. The dashes (-) in the formula are likewise not mathematical operators, but spacers, meaning "to": for instance the human formula is meaning that people may have 2 or 3 molars on each side of each jaw. 'd' denotes deciduous teeth (i.e. milk or baby teeth); lower case also indicates temporary teeth. Another annotation is , if the fact that it pertains to deciduous teeth is clearly stated, per examples found in some texts such as The Cambridge Dictionary of Human Biology and Evolution.
Permanent: This can also be written as . When the upper and lower dental formulae are the same, some texts write the formula without a fraction (in this case, 2.1.2.3), on the implicit assumption that the reader will realise it must apply to both upper and lower quadrants. This is seen, for example, throughout The Cambridge Dictionary of Human Biology and Evolution.
The greatest number of teeth in any known placental land mammal was 48, with a formula of . However, no living placental mammal has this number. In extant placental mammals, the maximum dental formula is for pigs. Mammalian tooth counts are usually identical in the upper and lower jaws, but not always. For example, the aye-aye has a formula of , demonstrating the need for both upper and lower quadrant counts.
Tooth naming discrepancies
Teeth are numbered starting at 1 in each group. Thus the human teeth are I1, I2, C1, P3, P4, M1, M2, and M3. (See next paragraph for premolar naming etymology.) In humans, the third molar is known as the wisdom tooth, whether or not it has erupted.
Regarding premolars, there is disagreement regarding whether the third type of deciduous tooth is a premolar (the general consensus among mammalogists) or a molar (commonly held among human anatomists). There is thus some discrepancy between nomenclature in zoology and in dentistry. This is because the terms of human dentistry, which have generally prevailed over time, have not included mammalian dental evolutionary theory. There were originally four premolars in each quadrant of early mammalian jaws. However, all living primates have lost at least the first premolar. "Hence most of the prosimians and platyrrhines have three premolars. Some genera have also lost more than one. A second premolar has been lost in all catarrhines. The remaining permanent premolars are then properly identified as P2, P3 and P4 or P3 and P4; however, traditional dentistry refers to them as P1 and P2".
Dental eruption sequence
The order in which teeth emerge through the gums is known as the dental eruption sequence. Rapidly developing anthropoid primates such as macaques, chimpanzees, and australopithecines have an eruption sequence of M1 I1 I2 M2 P3 P4 C M3, whereas anatomically modern humans have the sequence M1 I1 I2 C P3 P4 M2 M3. The later that tooth emergence begins, the earlier the anterior teeth (I1–P4) appear in the sequence.
Dental formulae examples
Dentition use in archaeology
Dentition, or the study of teeth, is an important area of study for archaeologists, especially those specializing in the study of older remains. Dentition affords many advantages over studying the rest of the skeleton itself (osteometry). The structure and arrangement of teeth is constant and, although it is inherited, does not undergo extensive change during environmental change, dietary specializations, or alterations in use patterns. The rest of the skeleton is much more likely to exhibit change because of adaptation. Teeth also preserve better than bone, and so the sample of teeth available to archaeologists is much more extensive and therefore more representative.
Dentition is particularly useful in tracking ancient populations' movements, because there are differences in the shapes of incisors, the number of grooves on molars, presence/absence of wisdom teeth, and extra cusps on particular teeth. These differences can not only be associated with different populations across space, but also change over time so that the study of the characteristics of teeth could say which population one is dealing with, and at what point in that population's history they are.
Dinosaurs
A dinosaur's dentition included all the teeth in its jawbones, which consist of the dentary, maxillary, and in some cases the premaxillary bones. The maxilla is the main bone of the upper jaw. The premaxilla is a smaller bone forming the anterior of the animal's upper jaw. The dentary is the main bone that forms the lower jaw (mandible). The predentary is a smaller bone that forms the anterior end of the lower jaw in ornithischian dinosaurs; it is always edentulous and supported a horny beak.
Unlike modern lizards, dinosaur teeth grew individually in the sockets of the jawbones, which are known as the dental alveoli. This thecodont dentition is also present in crocodilians and mammals, but is not found among the non-archosaur reptiles, which instead have acrodont or pleurodont dentition. Teeth that were lost were replaced by teeth below the roots in each tooth socket. Occlusion refers to the closing of the dinosaur's mouth, where the teeth from the upper and lower parts of the jaw meet. If the occlusion causes teeth from the maxillary or premaxillary bones to cover the teeth of the dentary and predentary, the dinosaur is said to have an overbite, the most common condition in this group. The opposite condition is considered to be an underbite, which is rare in theropod dinosaurs.
The majority of dinosaurs had teeth that were similarly shaped throughout their jaws but varied in size. Dinosaur tooth shapes included cylindrical, peg-like, teardrop-shaped, leaf-like, diamond-shaped and blade-like. A dinosaur that has a variety of tooth shapes is said to have heterodont dentition. An example of this are dinosaurs of the group Heterodontosauridae and the enigmatic early dinosaur, Eoraptor. While most dinosaurs had a single row of teeth on each side of their jaws, others had dental batteries where teeth in the cheek region were fused together to form compound teeth. Individually these teeth were not suitable for grinding food, but when joined together with other teeth they would form a large surface area for the mechanical digestion of tough plant materials. This type of dental strategy is observed in ornithopod and ceratopsian dinosaurs as well as the duck-billed hadrosaurs, which had more than one hundred teeth in each dental battery. The teeth of carnivorous dinosaurs, called ziphodont, were typically blade-like or cone-shaped, curved, with serrated edges. This dentition was adapted for grasping and cutting through flesh. In some cases, as observed in the railroad-spike-sized teeth of Tyrannosaurus rex, the teeth were designed to puncture and crush bone. Some dinosaurs had procumbent teeth, which projected forward in the mouth.
| Biology and health sciences | Gastrointestinal tract | Biology |
101336 | https://en.wikipedia.org/wiki/High-temperature%20superconductivity | High-temperature superconductivity | High-temperature superconductivity (high-c or HTS) is superconductivity in materials with a critical temperature (the temperature below which the material behaves as a superconductor) above , the boiling point of liquid nitrogen. They are only "high-temperature" relative to previously known superconductors, which function at colder temperatures, close to absolute zero. The "high temperatures" are still far below ambient (room temperature), and therefore require cooling. The first breakthrough of high-temperature superconductor was discovered in 1986 by IBM researchers Georg Bednorz and K. Alex Müller. Although the critical temperature is around , this new type of superconductor was readily modified by Ching-Wu Chu to make the first high-temperature superconductor with critical temperature . Bednorz and Müller were awarded the Nobel Prize in Physics in 1987 "for their important break-through in the discovery of superconductivity in ceramic materials". Most high-c materials are type-II superconductors.
The major advantage of high-temperature superconductors is that they can be cooled using liquid nitrogen, in contrast to the previously known superconductors that require expensive and hard-to-handle coolants, primarily liquid helium. A second advantage of high-c materials is they retain their superconductivity in higher magnetic fields than previous materials. This is important when constructing superconducting magnets, a primary application of high-c materials.
The majority of high-temperature superconductors are ceramic materials, rather than the previously known metallic materials. Ceramic superconductors are suitable for some practical uses but they still have many manufacturing issues. For example, most ceramics are brittle, which makes the fabrication of wires from them very problematic. However, overcoming these drawbacks is the subject of considerable research, and progress is ongoing.
The main class of high-temperature superconductors is copper oxides combined with other metals, especially the rare-earth barium copper oxides (REBCOs) such as yttrium barium copper oxide (YBCO). The second class of high-temperature superconductors in the practical classification is the iron-based compounds.
Magnesium diboride is sometimes included in high-temperature superconductors: It is relatively simple to manufacture, but it superconducts only below , which makes it unsuitable for liquid nitrogen cooling.
History
Superconductivity was discovered by Kamerlingh Onnes in 1911, in a metal solid. Ever since, researchers have attempted to observe superconductivity at increasing temperatures with the goal of finding a room-temperature superconductor. By the late 1970s, superconductivity was observed in several metallic compounds (in particular Nb-based, such as NbTi, Nb3Sn, and Nb3Ge) at temperatures that were much higher than those for elemental metals and which could even exceed .
In 1986, at the IBM research lab near Zürich in Switzerland, Bednorz and Müller were looking for superconductivity in a new class of ceramics: the copper oxides, or cuprates.
Bednorz encountered a particular copper oxide whose resistance dropped to zero at a temperature around . Their results were soon confirmed by many groups, notably Paul Chu at the University of Houston and Shoji Tanaka at the University of Tokyo.
In 1987, Philip W. Anderson gave the first theoretical description of these materials, based on the resonating valence bond (RVB) theory, but a full understanding of these materials is still developing today. These superconductors are now known to possess a d-wave pair symmetry. The first proposal that high-temperature cuprate superconductivity involves d-wave pairing was made in 1987 by N. E. Bickers, Douglas James Scalapino and R. T. Scalettar, followed by three subsequent theories in 1988 by Masahiko Inui, Sebastian Doniach, Peter J. Hirschfeld and Andrei E. Ruckenstein, using spin-fluctuation theory, and by Claudius Gros, Didier Poilblanc, Maurice T. Rice and FC. Zhang, and by Gabriel Kotliar and Jialin Liu identifying d-wave pairing as a natural consequence of the RVB theory. The confirmation of the d-wave nature of the cuprate superconductors was made by a variety of experiments, including the direct observation of the d-wave nodes in the excitation spectrum through angle resolved photoemission spectroscopy (ARPES), the observation of a half-integer flux in tunneling experiments, and indirectly from the temperature dependence of the penetration depth, specific heat and thermal conductivity.
As of 2021, the superconductor with the highest transition temperature at ambient pressure is the cuprate of mercury, barium, and calcium, at around . There are other superconductors with higher recorded transition temperaturesfor example lanthanum superhydride at , but these only occur at very high pressures.
The origin of high-temperature superconductivity is still not clear, but it seems that instead of electron–phonon attraction mechanisms, as in conventional superconductivity, one is dealing with genuine electronic mechanisms (e.g. by antiferromagnetic correlations), and instead of conventional, purely s-wave pairing, more exotic pairing symmetries are thought to be involved (d-wave in the case of the cuprates; primarily extended s-wave, but occasionally d-wave, in the case of the iron-based superconductors).
In 2014, evidence showing that fractional particles can happen in quasi two-dimensional magnetic materials, was found by École Polytechnique Fédérale de Lausanne (EPFL) scientists lending support for Anderson's theory of high-temperature superconductivity.
Selected list of superconductors
Properties
The "high-temperature" superconductor class has had many definitions.
The label high-c should be reserved for materials with critical temperatures greater than the boiling point of liquid nitrogen. However, a number of materialsincluding the original discovery and recently discovered pnictide superconductorshave critical temperatures below but nonetheless are commonly referred to in publications as high-c class.
A substance with a critical temperature above the boiling point of liquid nitrogen, together with a high critical magnetic field and critical current density (above which superconductivity is destroyed), would greatly benefit technological applications. In magnet applications, the high critical magnetic field may prove more valuable than the high c itself. Some cuprates have an upper critical field of about 100 tesla. However, cuprate materials are brittle ceramics that are expensive to manufacture and not easily turned into wires or other useful shapes. Furthermore, high-temperature superconductors do not form large, continuous superconducting domains, rather clusters of microdomains within which superconductivity occurs. They are therefore unsuitable for applications requiring actual superconductive currents, such as magnets for magnetic resonance spectrometers. For a solution to this (powders), see HTS wire.
There has been considerable debate regarding high-temperature superconductivity coexisting with magnetic ordering in YBCO, iron-based superconductors, several ruthenocuprates and other exotic superconductors, and the search continues for other families of materials. HTS are Type-II superconductors, which allow magnetic fields to penetrate their interior in quantized units of flux, meaning that much higher magnetic fields are required to suppress superconductivity. The layered structure also gives a directional dependence to the magnetic field response.
All known high-c superconductors are Type-II superconductors. In contrast to Type-I superconductors, which expel all magnetic fields due to the Meissner effect, Type-II superconductors allow magnetic fields to penetrate their interior in quantized units of flux, creating "holes" or "tubes" of normal metallic regions in the superconducting bulk called vortices. Consequently, high-c superconductors can sustain much higher magnetic fields.
Cuprates
Cuprates are layered materials, consisting of superconducting layers of copper oxide, separated by spacer layers.
Cuprates generally have a structure close to that of a two-dimensional material. Their superconducting properties are determined by electrons moving within weakly coupled copper-oxide (CuO2) layers. Neighbouring layers contain ions such as lanthanum, barium, strontium, or other atoms which act to stabilize the structures and dope electrons or holes onto the copper-oxide layers. The undoped "parent" or "mother" compounds are Mott insulators with long-range antiferromagnetic order at sufficiently low temperatures. Single band models are generally considered to be enough to describe the electronic properties.
The cuprate superconductors adopt a perovskite structure. The copper-oxide planes are checkerboard lattices with squares of O2− ions with a Cu2+ ion at the centre of each square. The unit cell is rotated by 45° from these squares. Chemical formulae of superconducting materials generally contain fractional numbers to describe the doping required for superconductivity. There are several families of cuprate superconductors and they can be categorized by the elements they contain and the number of adjacent copper-oxide layers in each superconducting block. For example, YBCO and BSCCO can alternatively be referred to as "Y123" and Bi2201/Bi2212/Bi2223 depending on the number of layers in each superconducting block (). The superconducting transition temperature has been found to peak at an optimal doping value (=0.16) and an optimal number of layers in each superconducting block, typically =3.
Possible mechanisms for superconductivity in the cuprates continue to be the subject of considerable debate and further research. Certain aspects common to all materials have been identified. Similarities between the antiferromagnetic the low-temperature state of undoped materials and the superconducting state that emerges upon doping, primarily the x2−y2 orbital state of the Cu2+ ions, suggest that electron–electron interactions are more significant than electron–phonon interactions in cupratesmaking the superconductivity unconventional. Recent work on the Fermi surface has shown that nesting occurs at four points in the antiferromagnetic Brillouin zone where spin waves exist and that the superconducting energy gap is larger at these points. The weak isotope effects observed for most cuprates contrast with conventional superconductors that are well described by BCS theory.
Similarities and differences in the properties of hole-doped and electron doped cuprates:
Presence of a pseudogap phase up to at least optimal doping.
Different trends in the Uemura plot relating transition temperature to the superfluid density. The inverse square of the London penetration depth appears to be proportional to the critical temperature for a large number of underdoped cuprate superconductors, but the constant of proportionality is different for hole- and electron-doped cuprates. The linear trend implies that the physics of these materials is strongly two-dimensional.
Universal hourglass-shaped feature in the spin excitations of cuprates measured using inelastic neutron diffraction.
Nernst effect evident in both the superconducting and pseudogap phases.
The electronic structure of superconducting cuprates is highly anisotropic (see the crystal structure of YBCO or BSCCO). Therefore, the Fermi surface of HTSC is very close to the Fermi surface of the doped CuO2 plane (or multi-planes, in case of multi-layer cuprates) and can be presented on the 2‑D reciprocal space (or momentum space) of the CuO2 lattice. The typical Fermi surface within the first CuO2 Brillouin zone is sketched in Fig. 1 (left). It can be derived from the band structure calculations or measured by angle resolved photoemission spectroscopy (ARPES). Fig. 1 (right) shows the Fermi surface of BSCCO measured by ARPES. In a wide range of charge carrier concentration (doping level), in which the hole-doped HTSC are superconducting, the Fermi surface is hole-like (i.e. open, as shown in Fig. 1). This results in an inherent in-plane anisotropy of the electronic properties of HTSC. In 2018, the full three dimensional Fermi surface structure was derived from soft x-ray ARPES.
Iron-based
Iron-based superconductors contain layers of iron and a pnictogensuch as arsenic or phosphorus, a chalcogen, or a crystallogen. This is currently the family with the second highest critical temperature, behind the cuprates. Interest in their superconducting properties began in 2006 with the discovery of superconductivity in LaFePO at and gained much greater attention in 2008 after the analogous material LaFeAs(O,F) was found to superconduct at up to under pressure.
The highest critical temperatures in the iron-based superconductor family exist in thin films of FeSe, where a critical temperature in excess of was reported in 2014.
Since the original discoveries several families of iron-based superconductors have emerged:
LnFeAs(O,F) or LnFeAsO1−x (Ln=lanthanide) with c up to , referred to as 1111 materials. A fluoride variant of these materials was subsequently found with similar c values.
(Ba,K)Fe2As2 and related materials with pairs of iron-arsenide layers, referred to as 122 compounds. c values range up to . These materials also superconduct when iron is replaced with cobalt.
LiFeAs and NaFeAs with c up to around . These materials superconduct close to stoichiometric composition and are referred to as 111 compounds.
FeSe with small off-stoichiometry or tellurium doping.
LaFeSiH with c around in its stoichiometric composition. This superconducting crystallogenide has oxide and fluoride variants LaFeSiOx and LaFeSiFx.
Most undoped iron-based superconductors show a tetragonal-orthorhombic structural phase transition followed at lower temperature by magnetic ordering, similar to the cuprate superconductors. However, they are poor metals rather than Mott insulators and have five bands at the Fermi surface rather than one.
The phase diagram emerging as the iron-arsenide layers are doped is remarkably similar, with the superconducting phase close to or overlapping the magnetic phase. Strong evidence that the c value varies with the As–Fe–As bond angles has already emerged and shows that the optimal c value is obtained with undistorted FeAs4 tetrahedra. The symmetry of the pairing wavefunction is still widely debated, but an extended s-wave scenario is currently favoured.
Magnesium diboride
Magnesium diboride is occasionally referred to as a high-temperature superconductor because its c value of is above that historically expected for BCS superconductors. However, it is more generally regarded as the highest c conventional superconductor, the increased c resulting from two separate bands being present at the Fermi level.
Carbon-based
In 1991 Hebard et al. discovered Fulleride superconductors, where alkali-metal atoms are intercalated into C60 molecules.
In 2008 Ganin et al. demonstrated superconductivity at temperatures of up to for Cs3C60.
P-doped Graphane was proposed in 2010 to be capable of sustaining high-temperature superconductivity.
On 31st of December 2023 "Global Room-Temperature Superconductivity in Graphite" was published in the journal "Advanced Quantum Technologies" claiming to demonstrate superconductivity at room temperature and ambient pressure in Highly oriented pyrolytic graphite with dense arrays of nearly parallel line defects.
Nickelates
In 1999, Anisimov et al. conjectured superconductivity in nickelates, proposing nickel oxides as direct analogs to the cuprate superconductors. Superconductivity in an infinite-layer nickelate, Nd0.8Sr0.2NiO2, was reported at the end of 2019 with a superconducting transition temperature between . This superconducting phase is observed in oxygen-reduced thin films created by the pulsed laser deposition of Nd0.8Sr0.2NiO3 onto SrTiO3 substrates that is then reduced to Nd0.8Sr0.2NiO2 via annealing the thin films at in the presence of CaH2. The superconducting phase is only observed in the oxygen reduced film and is not seen in oxygen reduced bulk material of the same stoichiometry, suggesting that the strain induced by the oxygen reduction of the Nd0.8Sr0.2NiO2 thin film changes the phase space to allow for superconductivity.
Of important is further to extract access hydrogen from the reduction with CaH2, otherwise topotactic hydrogen may prevent superconductivity.
Cuprates
The structure of cuprates which are superconductors are often closely related to perovskite structure, and the structure of these compounds has been described as a distorted, oxygen deficient multi-layered perovskite structure. One of the properties of the crystal structure of oxide superconductors is an alternating multi-layer of CuO2 planes with superconductivity taking place between these layers. The more layers of CuO2, the higher c. This structure causes a large anisotropy in normal conducting and superconducting properties, since electrical currents are carried by holes induced in the oxygen sites of the CuO2 sheets. The electrical conduction is highly anisotropic, with a much higher conductivity parallel to the CuO2 plane than in the perpendicular direction. Generally, critical temperatures depend on the chemical compositions, cations substitutions and oxygen content. They can be classified as superstripes; i.e., particular realizations of superlattices at atomic limit made of superconducting atomic layers, wires, dots separated by spacer layers, that gives multiband and multigap superconductivity.
Yttrium–barium cuprate
An yttrium–barium cuprate, YBa2Cu3O7−x (or Y123), was the first superconductor found above liquid nitrogen boiling point. There are two atoms of Barium for each atom of Yttrium.
The proportions of the three different metals in the YBa2Cu3O7 superconductor are in the mole ratio of 1 to 2 to 3 for yttrium to barium to copper, respectively: this particular superconductor has also often been referred to as the 123 superconductor.
The unit cell of YBa2Cu3O7 consists of three perovskite unit cells, which is pseudocubic, nearly orthorhombic. The other superconducting cuprates have another structure: they have a tetragonal cell.
Each perovskite cell contains a Y or Ba atom at the center: Ba in the bottom unit cell, Y in the middle one, and Ba in the top unit cell. Thus, Y and Ba are stacked in the sequence [Ba–Y–Ba] along the c-axis. All corner sites of the unit cell are occupied by Cu, which has two different coordinations, Cu(1) and Cu(2), with respect to oxygen. There are four possible crystallographic sites for oxygen: O(1), O(2), O(3) and O(4). The coordination polyhedra of Y and Ba with respect to oxygen are different. The tripling of the perovskite unit cell leads to nine oxygen atoms, whereas YBa2Cu3O7 has seven oxygen atoms and, therefore, is referred to as an oxygen-deficient perovskite structure. The structure has a stacking of different layers: (CuO)(BaO)(CuO2)(Y)(CuO2)(BaO)(CuO). One of the key feature of the unit cell of YBa2Cu3O7−x (YBCO) is the presence of two layers of CuO2. The role of the Y plane is to serve as a spacer between two CuO2 planes. In YBCO, the Cu–O chains are known to play an important role for superconductivity. c is maximal near when x ≈ 0.15 and the structure is orthorhombic. Superconductivity disappears at x ≈ 0.6, where the structural transformation of YBCO occurs from orthorhombic to tetragonal.
Other cuprates
The preparation of other cuprates is more difficult than the YBCO preparation.
They also have a different crystal structure: they are tetragonal where YBCO is orthorhombic.
Problems in these superconductors arise because of the existence of three or more phases having a similar layered structure.
Moreover, the crystal structure of other tested cuprate superconductors are very similar. Like YBCO, the perovskite-type feature and the presence of simple copper oxide (CuO2) layers also exist in these superconductors. However, unlike YBCO, Cu–O chains are not present in these superconductors. The YBCO superconductor has an orthorhombic structure, whereas the other high-c superconductors have a tetragonal structure.
There are three main classes of superconducting cuprates: bismuth-based, thallium-based and mercury-based.
The second cuprate by practical importance is currently BSCCO, a compound of Bi–Sr–Ca–Cu–O. The content of bismuth and strontium creates some chemical issues.
It has three superconducting phases forming a homologous series as Bi2Sr2Can−1CunO4+2n+x (n=1, 2 and 3).
These three phases are Bi-2201, Bi-2212 and Bi-2223, having transition temperatures of , and , respectively, where the numbering system represent number of atoms for Bi Sr, Ca and Cu respectively. The two phases have a tetragonal structure which consists of two sheared crystallographic unit cells. The unit cell of these phases has double Bi–O planes which are stacked in a way that the Bi atom of one plane sits below the oxygen atom of the next consecutive plane. The Ca atom forms a layer within the interior of the CuO2 layers in both Bi-2212 and Bi-2223; there is no Ca layer in the Bi-2201 phase. The three phases differ with each other in the number of cuprate planes; Bi-2201, Bi-2212 and Bi-2223 phases have one, two and three CuO2 planes, respectively. The c axis lattice constants of these phases increases with the number of cuprate planes (see table below). The coordination of the Cu atom is different in the three phases. The Cu atom forms an octahedral coordination with respect to oxygen atoms in the 2201 phase, whereas in 2212, the Cu atom is surrounded by five oxygen atoms in a pyramidal arrangement. In the 2223 structure, Cu has two coordinations with respect to oxygen: one Cu atom is bonded with four oxygen atoms in square planar configuration and another Cu atom is coordinated with five oxygen atoms in a pyramidal arrangement.
Cuprate of Tl–Ba–Ca: The first series of the Tl-based superconductor containing one Tl–O layer has the general formula TlBa2Can−1CunO2n+3, whereas the second series containing two Tl–O layers has a formula of Tl2Ba2Can−1CunO2n+4 with n =1, 2 and 3. In the structure of Tl2Ba2CuO6 (Tl-2201), there is one CuO2 layer with the stacking sequence (Tl–O) (Tl–O) (Ba–O) (Cu–O) (Ba–O) (Tl–O) (Tl–O). In Tl2Ba2CaCu2O8 (Tl-2212), there are two Cu–O layers with a Ca layer in between. Similar to the Tl2Ba2CuO6 structure, Tl–O layers are present outside the Ba–O layers. In Tl2Ba2Ca2Cu3O10 (Tl-2223), there are three CuO2 layers enclosing Ca layers between each of these. In Tl-based superconductors, c is found to increase with the increase in CuO2 layers. However, the value of c decreases after four CuO2 layers in TlBa2Can−1CunO2n+3, and in the Tl2Ba2Can−1CunO2n+4 compound, it decreases after three CuO2 layers.
Cuprate of Hg–Ba–Ca The crystal structure of HgBa2CuO4 (Hg-1201), HgBa2CaCu2O6 (Hg-1212) and HgBa2Ca2Cu3O8 (Hg-1223) is similar to that of Tl-1201, Tl-1212 and Tl-1223, with Hg in place of Tl. It is noteworthy that the c of the Hg compound (Hg-1201) containing one CuO2 layer is much larger as compared to the one-CuO2-layer compound of thallium (Tl-1201). In the Hg-based superconductor, c is also found to increase as the CuO2 layer increases. For Hg-1201, Hg-1212 and Hg-1223, the values of c are 94, 128, and the record value at ambient pressure , respectively, as shown in table below. The observation that the c of Hg-1223 increases to under high pressure indicates that the c of this compound is very sensitive to the structure of the compound.
Preparation and manufacturing
The simplest method for preparing ceramic superconductors is a solid-state thermochemical reaction involving mixing, calcination and sintering.
The appropriate amounts of precursor powders, usually oxides and carbonates, are mixed thoroughly using a Ball mill. Solution chemistry processes such as coprecipitation, freeze-drying and sol–gel methods are alternative ways for preparing a homogeneous mixture. These powders are calcined in the temperature range from for several hours. The powders are cooled, reground and calcined again. This process is repeated several times to get homogeneous material. The powders are subsequently compacted to pellets and sintered. The sintering environment such as temperature, annealing time, atmosphere and cooling rate play a very important role in getting good high-c superconducting materials. The YBa2Cu3O7−x compound is prepared by calcination and sintering of a homogeneous mixture of Y2O3, BaCO3 and CuO in the appropriate atomic ratio. Calcination is done at , whereas sintering is done at in an oxygen atmosphere. The oxygen stoichiometry in this material is very crucial for obtaining a superconducting YBa2Cu3O7−x compound. At the time of sintering, the semiconducting tetragonal YBa2Cu3O6 compound is formed, which, on slow cooling in oxygen atmosphere, turns into superconducting YBa2Cu3O7−x. The uptake and loss of oxygen are reversible in YBa2Cu3O7−x. A fully oxygenated orthorhombic YBa2Cu3O7−x sample can be transformed into tetragonal YBa2Cu3O6 by heating in a vacuum at temperature above .
The preparation of Bi-, Tl- and Hg-based high-c superconductors is more difficult than the YBCO preparation. Problems in these superconductors arise because of the existence of three or more phases having a similar layered structure. Thus, syntactic intergrowth and defects such as stacking faults occur during synthesis and it becomes difficult to isolate a single superconducting phase. For Bi–Sr–Ca–Cu–O, it is relatively simple to prepare the Bi-2212 (c ≈ 85 K) phase, whereas it is very difficult to prepare a single phase of Bi-2223 (c ≈ 110 K). The Bi-2212 phase appears only after few hours of sintering at , but the larger fraction of the Bi-2223 phase is formed after a long reaction time of more than a week at . Although the substitution of Pb in the Bi–Sr–Ca–Cu–O compound has been found to promote the growth of the high-c phase, a long sintering time is still required.
Ongoing research
The question of how superconductivity arises in high-temperature superconductors is one of the major unsolved problems of theoretical condensed matter physics. The mechanism that causes the electrons in these crystals to form pairs is not known. Despite intensive research and many promising leads, an explanation has so far eluded scientists. One reason for this is that the materials in question are generally very complex, multi-layered crystals (for example, BSCCO), making theoretical modelling difficult.
Improving the quality and variety of samples also gives rise to considerable research, both with the aim of improved characterisation of the physical properties of existing compounds, and synthesizing new materials, often with the hope of increasing c. Technological research focuses on making HTS materials in sufficient quantities to make their use economically viable as well as in optimizing their properties in relation to applications.
Metallic hydrogen has been proposed as a room-temperature superconductor, some experimental observations have detected the occurrence of the Meissner effect. LK-99, copper-doped lead-apatite, has also been proposed as a room-temperature superconductor.
Theoretical models
There have been two representative theories for high-temperature or unconventional superconductivity.
Firstly, weak coupling theory suggests superconductivity emerges from antiferromagnetic spin fluctuations in a doped system.
According to this theory, the pairing wave function of the cuprate HTS should have a dx2-y2 symmetry. Thus, determining whether the pairing wave function has d-wave symmetry is essential to test the spin fluctuation mechanism. That is, if the HTS order parameter (a pairing wave function like in Ginzburg–Landau theory) does not have d-wave symmetry, then a pairing mechanism related to spin fluctuations can be ruled out. (Similar arguments can be made for iron-based superconductors but the different material properties allow a different pairing symmetry.) Secondly, there was the interlayer coupling model, according to which a layered structure consisting of BCS-type (s-wave symmetry) superconductors can enhance the superconductivity by itself. By introducing an additional tunnelling interaction between each layer, this model successfully explained the anisotropic symmetry of the order parameter as well as the emergence of the HTS. Thus, in order to solve this unsettled problem, there have been numerous experiments such as photoemission spectroscopy, NMR, specific heat measurements, etc. Up to date the results were ambiguous, some reports supported the d symmetry for the HTS whereas others supported the s symmetry. This muddy situation possibly originated from the indirect nature of the experimental evidence, as well as experimental issues such as sample quality, impurity scattering, twinning, etc.
This summary makes an implicit assumption: superconductive properties can be treated by mean-field theory. It also fails to mention that in addition to the superconductive gap, there is a second gap, the pseudogap. The cuprate layers are insulating, and the superconductors are doped with interlayer impurities to make them metallic. The superconductive transition temperature can be maximized by varying the dopant concentration. The simplest example is La2CuO4, which consist of alternating CuO2 and LaO layers which are insulating when pure. When 8% of the La is replaced by Sr, the latter act as dopants, contributing holes to the CuO2 layers, and making the sample metallic. The Sr impurities also act as electronic bridges, enabling interlayer coupling. Proceeding from this picture, some theories argue that the basic pairing interaction is still interaction with phonons, as in the conventional superconductors with Cooper pairs. While the undoped materials are antiferromagnetic, even a few percent of impurity dopants introduce a smaller pseudogap in the CuO2 planes which is also caused by phonons. The gap decreases with increasing charge carriers, and as it nears the superconductive gap, the latter reaches its maximum. The reason for the high transition temperature is then argued to be due to the percolating behaviour of the carriersthe carriers follow zig-zag percolative paths, largely in metallic domains in the CuO2 planes, until blocked by charge density wave domain walls, where they use dopant bridges to cross over to a metallic domain of an adjacent CuO2 plane. The transition temperature maxima are reached when the host lattice has weak bond-bending forces, which produce strong electron–phonon interactions at the interlayer dopants.
D symmetry in YBCO
An experiment based on flux quantization of a three-grain ring of YBa2Cu3O7 (YBCO) was proposed to test the symmetry of the order parameter in the HTS. The symmetry of the order parameter could best be probed at the junction interface as the Cooper pairs tunnel across a Josephson junction or weak link.
It was expected that a half-integer flux, that is, a spontaneous magnetization could only occur for a junction of d symmetry superconductors. But, even if the junction experiment is the strongest method to determine the symmetry of the HTS order parameter, the results have been ambiguous. John R. Kirtley and C. C. Tsuei thought that the ambiguous results came from the defects inside the HTS, so that they designed an experiment where both clean limit (no defects) and dirty limit (maximal defects) were considered simultaneously. In the experiment, the spontaneous magnetization was clearly observed in YBCO, which supported the d symmetry of the order parameter in YBCO. But, since YBCO is orthorhombic, it might inherently have an admixture of s symmetry. So, by tuning their technique further, they found that there was an admixture of s symmetry in YBCO within about 3%. Also, they found that there was a pure dx2−y2 order parameter symmetry in the tetragonal Tl2Ba2CuO6.
Spin-fluctuation mechanism
Despite all these years, the mechanism of high-c superconductivity is still highly controversial, mostly due to the lack of exact theoretical computations on such strongly interacting electron systems. However, most rigorous theoretical calculations, including phenomenological and diagrammatic approaches, converge on magnetic fluctuations as the pairing mechanism for these systems. The qualitative explanation is as follows:
In a superconductor, the flow of electrons cannot be resolved into individual electrons, but instead consists of many pairs of bound electrons, called Cooper pairs. In conventional superconductors, these pairs are formed when an electron moving through the material distorts the surrounding crystal lattice, which in turn attracts another electron and forms a bound pair. This is sometimes called the "water bed" effect. Each Cooper pair requires a certain minimum energy to be displaced, and if the thermal fluctuations in the crystal lattice are smaller than this energy the pair can flow without dissipating energy. This ability of the electrons to flow without resistance leads to superconductivity.
In a high-c superconductor, the mechanism is extremely similar to a conventional superconductor, except, in this case, phonons virtually play no role and their role is replaced by spin-density waves. Just as all known conventional superconductors are strong phonon systems, all known high-c superconductors are strong spin-density wave systems, within close vicinity of a magnetic transition to, for example, an antiferromagnet. When an electron moves in a high-c superconductor, its spin creates a spin-density wave around it. This spin-density wave in turn causes a nearby electron to fall into the spin depression created by the first electron (water-bed effect again). Hence, again, a Cooper pair is formed. When the system temperature is lowered, more spin density waves and Cooper pairs are created, eventually leading to superconductivity. Note that in high-c systems, as these systems are magnetic systems due to the Coulomb interaction, there is a strong Coulomb repulsion between electrons. This Coulomb repulsion prevents pairing of the Cooper pairs on the same lattice site. The pairing of the electrons occur at near-neighbor lattice sites as a result. This is the so-called d-wave pairing, where the pairing state has a node (zero) at the origin.
Examples
Examples of high-c cuprate superconductors include YBCO and BSCCO, which are the most known materials that achieve superconductivity above the boiling point of liquid nitrogen.
| Physical sciences | Electrical circuits | Physics |
101453 | https://en.wikipedia.org/wiki/Dirichlet%27s%20theorem%20on%20arithmetic%20progressions | Dirichlet's theorem on arithmetic progressions | In number theory, Dirichlet's theorem, also called the Dirichlet prime number theorem, states that for any two positive coprime integers a and d, there are infinitely many primes of the form a + nd, where n is also a positive integer. In other words, there are infinitely many primes that are congruent to a modulo d. The numbers of the form a + nd form an arithmetic progression
and Dirichlet's theorem states that this sequence contains infinitely many prime numbers. The theorem extends Euclid's theorem that there are infinitely many prime numbers (of the form 1 + 2n). Stronger forms of Dirichlet's theorem state that for any such arithmetic progression, the sum of the reciprocals of the prime numbers in the progression diverges and that different such arithmetic progressions with the same modulus have approximately the same proportions of primes. Equivalently, the primes are evenly distributed (asymptotically) among the congruence classes modulo d containing a'''s coprime to d.
The theorem is named after the German mathematician Peter Gustav Lejeune Dirichlet, who proved it in 1837.
Examples
The primes of the form 4n + 3 are
3, 7, 11, 19, 23, 31, 43, 47, 59, 67, 71, 79, 83, 103, 107, 127, 131, 139, 151, 163, 167, 179, 191, 199, 211, 223, 227, 239, 251, 263, 271, 283, ...
They correspond to the following values of n:
0, 1, 2, 4, 5, 7, 10, 11, 14, 16, 17, 19, 20, 25, 26, 31, 32, 34, 37, 40, 41, 44, 47, 49, 52, 55, 56, 59, 62, 65, 67, 70, 76, 77, 82, 86, 89, 91, 94, 95, ...
The strong form of Dirichlet's theorem implies that
is a divergent series.
Sequences dn + a with odd d are often ignored because half the numbers are even and the other half is the same numbers as a sequence with 2d, if we start with n = 0. For example, 6n + 1 produces the same primes as 3n + 1, while 6n + 5 produces the same as 3n + 2 except for the only even prime 2. The following table lists several arithmetic progressions with infinitely many primes and the first few ones in each of them.
We can generate some forms of primes by using an iterative method. For example, we can generate primes of the form by using the following method:
Let . Then we let which is prime. We continue by computing . Because is of the form , either 13 or 67 is of the form . We have that and is prime, so . We then continue this process to find successive primes of the form (Silverman 2013).
Distribution
Since the primes thin out, on average, in accordance with the prime number theorem, the same must be true for the primes in arithmetic progressions. It is natural to ask about the way the primes are shared between the various arithmetic progressions for a given value of d (there are d of those, essentially, if we do not distinguish two progressions sharing almost all their terms). The answer is given in this form: the number of feasible progressions modulo d — those where a and d do not have a common factor > 1 — is given by Euler's totient function
Further, the proportion of primes in each of those is
For example, if d is a prime number q, each of the q − 1 progressions
(all except )
contains a proportion 1/(q − 1) of the primes.
When compared to each other, progressions with a quadratic nonresidue remainder have typically slightly more elements than those with a quadratic residue remainder (Chebyshev's bias).
History
In 1737, Euler related the study of prime numbers to what is known now as the Riemann zeta function: he showed that the value reduces to a ratio of two infinite products, Π p / Π (p–1), for all primes p, and that the ratio is infinite.Sandifer, C. Edward, The Early Mathematics of Leonhard Euler (Washington, D.C.: The Mathematical Association of America, 2007), p. 253. In 1775, Euler stated the theorem for the cases of a + nd, where a = 1.
This special case of Dirichlet's theorem can be proven using cyclotomic polynomials.
The general form of the theorem was first conjectured by Legendre in his attempted unsuccessful proofs of quadratic reciprocity — as Gauss noted in his Disquisitiones Arithmeticae — but it was proved by with Dirichlet L-series. The proof is modeled on Euler's earlier work relating the Riemann zeta function to the distribution of primes. The theorem represents the beginning of rigorous analytic number theory.
gave an elementary proof.
Proof
Dirichlet's theorem is proved by showing that the value of the Dirichlet L-function (of a non-trivial character) at 1 is nonzero. The proof of this statement requires some calculus and analytic number theory . The particular case a = 1 (i.e., concerning the primes that are congruent to 1 modulo some n) can be proven by analyzing the splitting behavior of primes in cyclotomic extensions, without making use of calculus .
Although the proof of Dirichlet's Theorem makes use of calculus and analytic number theory, some proofs of examples are much more straightforward. In particular, the proof of the example of infinitely many primes of the form makes an argument similar to the one made in the proof of Euclid's theorem (Silverman 2013). The proof is given below:
We want to prove that there are infinitely many primes of the form . Assume, for contradiction, that there are only finitely many primes of the form . We then compile a list of all such primes where . Let . It is clear that none of the primes in the list divide . Next, suppose that is composite. Then has unique prime factorization where each is prime. Because , is odd and must be the product of only odd primes. Any odd prime must be such that or . It cannot be that because if this were the case, then . So there exists a prime such that . However, none of the primes in the list divide , a contradiction. Therefore must be prime and . Hence, is a prime of the form , but it isn't included in the list . Thus, the list doesn't contain all such primes and there must be infinitely many primes of the form (Silverman 2013).
Generalizations
The Bunyakovsky conjecture generalizes Dirichlet's theorem to higher-degree polynomials. Whether or not even simple quadratic polynomials such as (known from Landau's fourth problem) attain infinitely many prime values is an important open problem.
Dickson's conjecture generalizes Dirichlet's theorem to more than one polynomial.
Schinzel's hypothesis H generalizes these two conjectures, i.e. generalizes to more than one polynomial with degree larger than one.
In algebraic number theory, Dirichlet's theorem generalizes to the Chebotarev's density theorem.
Linnik's theorem (1944) concerns the size of the smallest prime in a given arithmetic progression. Linnik proved that the progression a + nd (as n ranges through the positive integers) contains a prime of magnitude at most cdL for absolute constants c and L. Subsequent researchers have reduced L to 5.
An analogue of Dirichlet's theorem holds in the framework of dynamical systems (T. Sunada and A. Katsuda, 1990).
Shiu showed that any arithmetic progression satisfying the hypothesis of Dirichlet's theorem will in fact contain arbitrarily long runs of consecutive'' prime numbers.
| Mathematics | Sequences | null |
101863 | https://en.wikipedia.org/wiki/Linear%20independence | Linear independence | In the theory of vector spaces, a set of vectors is said to be if there exists no nontrivial linear combination of the vectors that equals the zero vector. If such a linear combination exists, then the vectors are said to be . These concepts are central to the definition of dimension.
A vector space can be of finite dimension or infinite dimension depending on the maximum number of linearly independent vectors. The definition of linear dependence and the ability to determine whether a subset of vectors in a vector space is linearly dependent are central to determining the dimension of a vector space.
Definition
A sequence of vectors from a vector space is said to be linearly dependent, if there exist scalars not all zero, such that
where denotes the zero vector.
This implies that at least one of the scalars is nonzero, say , and the above equation is able to be written as
if and if
Thus, a set of vectors is linearly dependent if and only if one of them is zero or a linear combination of the others.
A sequence of vectors is said to be linearly independent if it is not linearly dependent, that is, if the equation
can only be satisfied by for This implies that no vector in the sequence can be represented as a linear combination of the remaining vectors in the sequence. In other words, a sequence of vectors is linearly independent if the only representation of as a linear combination of its vectors is the trivial representation in which all the scalars are zero. Even more concisely, a sequence of vectors is linearly independent if and only if can be represented as a linear combination of its vectors in a unique way.
If a sequence of vectors contains the same vector twice, it is necessarily dependent. The linear dependency of a sequence of vectors does not depend of the order of the terms in the sequence. This allows defining linear independence for a finite set of vectors: A finite set of vectors is linearly independent if the sequence obtained by ordering them is linearly independent. In other words, one has the following result that is often useful.
A sequence of vectors is linearly independent if and only if it does not contain the same vector twice and the set of its vectors is linearly independent.
Infinite case
An infinite set of vectors is linearly independent if every nonempty finite subset is linearly independent. Conversely, an infinite set of vectors is linearly dependent if it contains a finite subset that is linearly dependent, or equivalently, if some vector in the set is a linear combination of other vectors in the set.
An indexed family of vectors is linearly independent if it does not contain the same vector twice, and if the set of its vectors is linearly independent. Otherwise, the family is said to be linearly dependent.
A set of vectors which is linearly independent and spans some vector space, forms a basis for that vector space. For example, the vector space of all polynomials in over the reals has the (infinite) subset as a basis.
Geometric examples
and are independent and define the plane P.
, and are dependent because all three are contained in the same plane.
and are dependent because they are parallel to each other.
, and are independent because and are independent of each other and is not a linear combination of them or, equivalently, because they do not belong to a common plane. The three vectors define a three-dimensional space.
The vectors (null vector, whose components are equal to zero) and are dependent since
Geographic location
A person describing the location of a certain place might say, "It is 3 miles north and 4 miles east of here." This is sufficient information to describe the location, because the geographic coordinate system may be considered as a 2-dimensional vector space (ignoring altitude and the curvature of the Earth's surface). The person might add, "The place is 5 miles northeast of here." This last statement is true, but it is not necessary to find the location.
In this example the "3 miles north" vector and the "4 miles east" vector are linearly independent. That is to say, the north vector cannot be described in terms of the east vector, and vice versa. The third "5 miles northeast" vector is a linear combination of the other two vectors, and it makes the set of vectors linearly dependent, that is, one of the three vectors is unnecessary to define a specific location on a plane.
Also note that if altitude is not ignored, it becomes necessary to add a third vector to the linearly independent set. In general, linearly independent vectors are required to describe all locations in -dimensional space.
Evaluating linear independence
The zero vector
If one or more vectors from a given sequence of vectors is the zero vector then the vector are necessarily linearly dependent (and consequently, they are not linearly independent).
To see why, suppose that is an index (i.e. an element of ) such that Then let (alternatively, letting be equal any other non-zero scalar will also work) and then let all other scalars be (explicitly, this means that for any index other than (i.e. for ), let so that consequently ).
Simplifying gives:
Because not all scalars are zero (in particular, ), this proves that the vectors are linearly dependent.
As a consequence, the zero vector can not possibly belong to any collection of vectors that is linearly independent.
Now consider the special case where the sequence of has length (i.e. the case where ).
A collection of vectors that consists of exactly one vector is linearly dependent if and only if that vector is zero.
Explicitly, if is any vector then the sequence (which is a sequence of length ) is linearly dependent if and only if alternatively, the collection is linearly independent if and only if
Linear dependence and independence of two vectors
This example considers the special case where there are exactly two vector and from some real or complex vector space.
The vectors and are linearly dependent if and only if at least one of the following is true:
is a scalar multiple of (explicitly, this means that there exists a scalar such that ) or
is a scalar multiple of (explicitly, this means that there exists a scalar such that ).
If then by setting we have (this equality holds no matter what the value of is), which shows that (1) is true in this particular case. Similarly, if then (2) is true because
If (for instance, if they are both equal to the zero vector ) then both (1) and (2) are true (by using for both).
If then is only possible if and ; in this case, it is possible to multiply both sides by to conclude
This shows that if and then (1) is true if and only if (2) is true; that is, in this particular case either both (1) and (2) are true (and the vectors are linearly dependent) or else both (1) and (2) are false (and the vectors are linearly independent).
If but instead then at least one of and must be zero.
Moreover, if exactly one of and is (while the other is non-zero) then exactly one of (1) and (2) is true (with the other being false).
The vectors and are linearly independent if and only if is not a scalar multiple of and is not a scalar multiple of .
Vectors in R2
Three vectors: Consider the set of vectors and then the condition for linear dependence seeks a set of non-zero scalars, such that
or
Row reduce this matrix equation by subtracting the first row from the second to obtain,
Continue the row reduction by (i) dividing the second row by 5, and then (ii) multiplying by 3 and adding to the first row, that is
Rearranging this equation allows us to obtain
which shows that non-zero ai exist such that can be defined in terms of and Thus, the three vectors are linearly dependent.
Two vectors: Now consider the linear dependence of the two vectors and and check,
or
The same row reduction presented above yields,
This shows that which means that the vectors and are linearly independent.
Vectors in R4
In order to determine if the three vectors in
are linearly dependent, form the matrix equation,
Row reduce this equation to obtain,
Rearrange to solve for v3 and obtain,
This equation is easily solved to define non-zero ai,
where can be chosen arbitrarily. Thus, the vectors and are linearly dependent.
Alternative method using determinants
An alternative method relies on the fact that vectors in are linearly independent if and only if the determinant of the matrix formed by taking the vectors as its columns is non-zero.
In this case, the matrix formed by the vectors is
We may write a linear combination of the columns as
We are interested in whether for some nonzero vector Λ. This depends on the determinant of , which is
Since the determinant is non-zero, the vectors and are linearly independent.
Otherwise, suppose we have vectors of coordinates, with Then A is an n×m matrix and Λ is a column vector with entries, and we are again interested in AΛ = 0. As we saw previously, this is equivalent to a list of equations. Consider the first rows of , the first equations; any solution of the full list of equations must also be true of the reduced list. In fact, if is any list of rows, then the equation must be true for those rows.
Furthermore, the reverse is true. That is, we can test whether the vectors are linearly dependent by testing whether
for all possible lists of rows. (In case , this requires only one determinant, as above. If , then it is a theorem that the vectors must be linearly dependent.) This fact is valuable for theory; in practical calculations more efficient methods are available.
More vectors than dimensions
If there are more vectors than dimensions, the vectors are linearly dependent. This is illustrated in the example above of three vectors in
Natural basis vectors
Let and consider the following elements in , known as the natural basis vectors:
Then are linearly independent.
Linear independence of functions
Let be the vector space of all differentiable functions of a real variable . Then the functions and in are linearly independent.
Proof
Suppose and are two real numbers such that
Take the first derivative of the above equation:
for values of We need to show that and In order to do this, we subtract the first equation from the second, giving . Since is not zero for some , It follows that too. Therefore, according to the definition of linear independence, and are linearly independent.
Space of linear dependencies
A linear dependency or linear relation among vectors is a tuple with scalar components such that
If such a linear dependence exists with at least a nonzero component, then the vectors are linearly dependent. Linear dependencies among form a vector space.
If the vectors are expressed by their coordinates, then the linear dependencies are the solutions of a homogeneous system of linear equations, with the coordinates of the vectors as coefficients. A basis of the vector space of linear dependencies can therefore be computed by Gaussian elimination.
Generalizations
Affine independence
A set of vectors is said to be affinely dependent if at least one of the vectors in the set can be defined as an affine combination of the others. Otherwise, the set is called affinely independent. Any affine combination is a linear combination; therefore every affinely dependent set is linearly dependent. Conversely, every linearly independent set is affinely independent.
Consider a set of vectors of size each, and consider the set of augmented vectors of size each. The original vectors are affinely independent if and only if the augmented vectors are linearly independent.
Linearly independent vector subspaces
Two vector subspaces and of a vector space are said to be if
More generally, a collection of subspaces of are said to be if for every index where
The vector space is said to be a of if these subspaces are linearly independent and
| Mathematics | Linear algebra | null |
101965 | https://en.wikipedia.org/wiki/Psoriasis | Psoriasis | Psoriasis is a long-lasting, noncontagious autoimmune disease characterized by patches of abnormal skin. These areas are red, pink, or purple, dry, itchy, and scaly. Psoriasis varies in severity from small localized patches to complete body coverage. Injury to the skin can trigger psoriatic skin changes at that spot, which is known as the Koebner phenomenon.
The five main types of psoriasis are plaque, guttate, inverse, pustular, and erythrodermic. Plaque psoriasis, also known as psoriasis vulgaris, makes up about 90% of cases. It typically presents as red patches with white scales on top. Areas of the body most commonly affected are the back of the forearms, shins, navel area, and scalp. Guttate psoriasis has drop-shaped lesions. Pustular psoriasis presents as small, noninfectious, pus-filled blisters. Inverse psoriasis forms red patches in skin folds. Erythrodermic psoriasis occurs when the rash becomes very widespread and can develop from any of the other types. Fingernails and toenails are affected in most people with psoriasis at some point in time. This may include pits in the nails or changes in nail color.
Psoriasis is generally thought to be a genetic disease that is triggered by environmental factors. If one twin has psoriasis, the other twin is three times more likely to be affected if the twins are identical than if they are nonidentical. This suggests that genetic factors predispose to psoriasis. Symptoms often worsen during winter and with certain medications, such as beta blockers or NSAIDs. Infections and psychological stress can also play a role. The underlying mechanism involves the immune system reacting to skin cells. Diagnosis is typically based on the signs and symptoms.
There is no known cure for psoriasis, but various treatments can help control the symptoms. These treatments include steroid creams, vitamin D3 cream, ultraviolet light, immunosuppressive drugs, such as methotrexate, and biologic therapies targeting specific immunologic pathways. About 75% of skin involvement improves with creams alone. The disease affects 2–4% of the population. Men and women are affected with equal frequency. The disease may begin at any age, but typically starts in adulthood. Psoriasis is associated with an increased risk of psoriatic arthritis, lymphomas, cardiovascular disease, Crohn's disease, and depression. Psoriatic arthritis affects up to 30% of individuals with psoriasis.
The word "psoriasis" is from Greek ψωρίασις, meaning "itching condition" or "being itchy" from psora, "itch", and -iasis, "action, condition".
Signs and symptoms
Plaque psoriasis
Psoriasis vulgaris (also known as chronic stationary psoriasis or plaque-like psoriasis) is the most common form and affects 85–90% of people with psoriasis. Plaque psoriasis typically appears as raised areas of inflamed skin covered with silvery-white, scaly skin. These areas are called plaques and are most commonly found on the elbows, knees, scalp, and back.
Other forms
Additional types of psoriasis comprise about 10% of cases. They include pustular, inverse, napkin, guttate, oral, and seborrheic-like forms.
Pustular psoriasis
Pustular psoriasis appears as raised bumps filled with noninfectious pus (pustules). The skin under and surrounding the pustules is red and tender. Pustular psoriasis can either be localized or more widespread throughout the body. Two types of localized pustular psoriasis include psoriasis pustulosa palmoplantaris and acrodermatitis continua of Hallopeau; both forms are localized to the hands and feet.
Inverse psoriasis
Inverse psoriasis (also known as flexural psoriasis) appears as smooth, inflamed patches of skin. The patches frequently affect skin folds, particularly around the genitals (between the thigh and groin), the armpits, in the skin folds of an overweight abdomen (known as panniculus), between the buttocks in the intergluteal cleft, and under the breasts in the inframammary fold. Heat, trauma, and infection are thought to play a role in the development of this atypical form of psoriasis.
Napkin psoriasis
Napkin psoriasis is a subtype of psoriasis common in infants under the age of two and is characterized by red papules with silver scales in the diaper area that may extend to the torso or limbs. Napkin psoriasis is often misdiagnosed as napkin dermatitis (diaper rash). It typically improves as children age and may later present in more common forms as plaque psoriasis or inverse psoriasis.
Guttate psoriasis
Guttate psoriasis is an inflammatory condition characterized by numerous small, scaly, red or pink, droplet-like lesions (papules). These numerous papules appear over large areas of the body, primarily the trunk, limbs, and scalp, but typically spare the palms and soles. Guttate psoriasis is often triggered by a streptococcal infection (oropharyngeal or perianal) and typically occurs 1–3 weeks post-infection. Guttate psoriasis is most commonly seen in children and young adults and diagnosis is typically made based on history and clinical exam findings. Skin biopsy can also be performed which typically shows a psoriasiform reaction pattern characterized by epidermal hyperplasia with elongation of the rete ridges.
There is no firm evidence regarding the best management for guttate psoriasis; however, first-line therapy for mild guttate psoriasis typically includes topical corticosteroids. Phototherapy can be used for moderate or severe guttate psoriasis. Biologic treatments have not been well studied in the treatment of guttate psoriasis.
Guttate psoriasis has a better prognosis than plaque psoriasis and typically resolves within 1–3 weeks; however, up to 40% of patients with guttate psoriasis eventually convert to plaque psoriasis.
Erythrodermic psoriasis
Psoriatic erythroderma (erythrodermic psoriasis) involves widespread inflammation and exfoliation of the skin over most of the body surface, often involving greater than 90% of the body surface area. It may be accompanied by severe dryness, itching, swelling, and pain. It can develop from any type of psoriasis. It is often the result of an exacerbation of unstable plaque psoriasis, particularly following the abrupt withdrawal of systemic glucocorticoids. This form of psoriasis can be fatal as the extreme inflammation and exfoliation disrupt the body's ability to regulate temperature and perform barrier functions.
Mouth
Psoriasis in the mouth is very rare, in contrast to lichen planus, another common papulosquamous disorder that commonly involves both the skin and mouth. When psoriasis involves the oral mucosa (the lining of the mouth), it may be asymptomatic, but it may appear as white or grey-yellow plaques. Fissured tongue is the most common finding in those with oral psoriasis and has been reported to occur in 6.5–20% of people with psoriasis affecting the skin. The microscopic appearance of oral mucosa affected by geographic tongue (migratory stomatitis) is very similar to the appearance of psoriasis. A recent study found an association between the two conditions, and it suggests that geographic tongue might be a predictor to psoriasis.
Seborrheic-like psoriasis
Seborrheic-like psoriasis is a common form of psoriasis with clinical aspects of psoriasis and seborrheic dermatitis, and it may be difficult to distinguish from the latter. This form of psoriasis typically manifests as red plaques with greasy scales in areas of higher sebum production such as the scalp, forehead, skin folds next to the nose, the skin surrounding the mouth, skin on the chest above the sternum, and in skin folds.
Psoriatic arthritis
Psoriatic arthritis is a form of chronic inflammatory arthritis that has a highly variable clinical presentation and frequently occurs in association with skin and nail psoriasis. It typically involves painful inflammation of the joints and surrounding connective tissue and can occur in any joint, but most commonly affects the joints of the fingers and toes. This can result in a sausage-shaped swelling of the fingers and toes known as dactylitis. Psoriatic arthritis can also affect the hips, knees, spine (spondylitis), and sacroiliac joint (sacroiliitis). About 30% of individuals with psoriasis will develop psoriatic arthritis. Skin manifestations of psoriasis tend to occur before arthritic manifestations in about 75% of cases.
Nail changes
Psoriasis can affect the nails and produces a variety of changes in the appearance of fingers and toenails. Nail psoriasis occurs in 40–45% of people with psoriasis affecting the skin, and has a lifetime incidence of 80–90% in those with psoriatic arthritis. These changes include pitting of the nails (pinhead-sized depressions in the nail is seen in 70% with nail psoriasis), whitening of the nail, small areas of bleeding from capillaries under the nail, yellow-reddish discoloration of the nails known as the oil drop or salmon spots, dryness, thickening of the skin under the nail (subungual hyperkeratosis), loosening and separation of the nail (onycholysis), and crumbling of the nail.
Medical signs
In addition to the appearance and distribution of the rash, specific medical signs may be used by medical practitioners to assist with diagnosis. These may include Auspitz's sign (pinpoint bleeding when the scale is removed), Koebner phenomenon (psoriatic skin lesions induced by trauma to the skin), and itching and pain localized to papules and plaques.
Causes
The cause of psoriasis is not fully understood. Genetics, seasonal changes, skin damage, climate, immunocompromised state, specific infections, and the use of some medications have been connected with different types of psoriasis.
Genetics
Around one-third of people with psoriasis report a family history of the disease, and researchers have identified genetic loci associated with the condition. Identical twin studies suggest a 70% chance of a twin developing psoriasis if the other twin has the disorder. The risk is around 20% for fraternal twins. These findings suggest both a genetic susceptibility and an environmental response in developing psoriasis.
Psoriasis has a strong hereditary component, and many genes are associated with it, but how those genes work together is unclear. Most of the identified genes relate to the immune system, particularly the major histocompatibility complex (MHC) and T cells. Genetic studies are valuable due to their ability to identify molecular mechanisms and pathways for further study and potential medication targets.
Classic genome-wide linkage analysis has identified nine loci on different chromosomes associated with psoriasis. They are called psoriasis susceptibility 1 through 9 (PSORS1 through PSORS9). Within those loci are genes on pathways that lead to inflammation. Certain variations (mutations) of those genes are commonly found in psoriasis. Genome-wide association scans have identified other genes that are altered to characteristic variants in psoriasis. Some of these genes express inflammatory signal proteins, which affect cells in the immune system that are also involved in psoriasis. Some of these genes are also involved in other autoimmune diseases.
The major determinant is PSORS1, which probably accounts for 35–50% of psoriasis heritability. It controls genes that affect the immune system or encode skin proteins that are overabundant with psoriasis. PSORS1 is located on chromosome 6 in the MHC, which controls important immune functions. Three genes in the PSORS1 locus have a strong association with psoriasis vulgaris: HLA-C variant HLA-Cw6, which encodes an MHC class I protein; CCHCR1, variant WWC, which encodes a coiled coil protein overexpressed in psoriatic epidermis; and CDSN, variant allele 5, which encodes corneodesmosin, a protein expressed in the granular and cornified layers of the epidermis and upregulated in psoriasis.
Two major immune system genes under investigation are interleukin-12 subunit beta (IL12B) on chromosome 5q, which expresses interleukin-12B; and IL23R on chromosome 1p, which expresses the interleukin-23 receptor and is involved in T cell differentiation. Interleukin-23 receptor and IL12B have both been strongly linked with psoriasis. T cells are involved in the inflammatory process that leads to psoriasis. These genes are on the pathway that upregulates tumor necrosis factor-α and nuclear factor κB, two genes involved in inflammation. The first gene directly linked to psoriasis was identified as the CARD14 gene located in the PSORS2 locus. A rare mutation in the gene encoding for the CARD14-regulated protein plus an environmental trigger was enough to cause plaque psoriasis (the most common form of psoriasis).
Lifestyle
Conditions reported as worsening the disease include chronic infections, stress, and changes in season and climate. Other factors that might worsen the condition include hot water, scratching psoriasis skin lesions, skin dryness, excessive alcohol consumption, cigarette smoking, and obesity. The effects of stopping cigarette smoking or alcohol misuse have yet to be studied as of 2019.
HIV
The rate of psoriasis in human immunodeficiency virus-positive (HIV) individuals is comparable to that of HIV-negative individuals, but psoriasis tends to be more severe in people infected with HIV. A much higher rate of psoriatic arthritis occurs in HIV-positive individuals with psoriasis than in those without the infection. The immune response in those infected with HIV is typically characterized by cellular signals from Th2 subset of CD4+ helper T cells, whereas the immune response in psoriasis vulgaris is characterized by a pattern of cellular signals typical of Th1 subset of CD4+ helper T cells and Th17 helper T cells. The diminished CD4+-T cell presence is thought to cause overactivation of CD8+-T cells, which are responsible for the exacerbation of psoriasis in HIV-positive people. Psoriasis in those with HIV/AIDS is often severe and may be untreatable with conventional therapy. In those with long-term, well-controlled psoriasis, new HIV infection can trigger a severe flare-up of psoriasis and/or psoriatic arthritis.
Microbes
Psoriasis has been described as occurring after strep throat, and may be worsened by skin or gut colonization with Staphylococcus aureus, Malassezia spp., and Candida albicans. Guttate psoriasis often affects children and adolescents and can be triggered by a recent group A streptococcal infection (tonsillitis or pharyngitis).
Medications
Drug-induced psoriasis may occur with beta blockers, lithium, antimalarial medications, nonsteroidal anti-inflammatory drugs, terbinafine, calcium channel blockers, captopril, glyburide, granulocyte colony-stimulating factor, interleukins, interferons, lipid-lowering medications, and paradoxically TNF inhibitors such as infliximab or adalimumab. Withdrawal of corticosteroids (topical steroid cream) can aggravate psoriasis due to the rebound effect.
Pathophysiology
Psoriasis is characterized by an abnormally excessive and rapid growth of the epidermal layer of the skin. Abnormal production of skin cells (especially during wound repair) and an overabundance of skin cells result from the sequence of pathological events in psoriasis. The sequence of pathological events in psoriasis is thought to start with an initiation phase in which an event (skin trauma, infection, or drugs) leads to activation of the immune system and then the maintenance phase consisting of chronic progression of the disease. Skin cells are replaced every 3–5 days in psoriasis rather than the usual 28–30 days. These changes are believed to stem from the premature maturation of keratinocytes induced by an inflammatory cascade in the dermis involving dendritic cells, macrophages, and T cells (three subtypes of immune cells). These immune cells move from the dermis to the epidermis and secrete inflammatory chemical signals (cytokines) such as interleukin-36γ, tumor necrosis factor-α, interleukin-1β, interleukin-6, and interleukin-22. These secreted inflammatory signals are believed to stimulate keratinocytes to proliferate. One hypothesis is that psoriasis involves a defect in regulatory T cells, and in the regulatory cytokine interleukin-10. The inflammatory cytokines found in psoriatic nails and joints (in the case of psoriatic arthritis) are similar to those of psoriatic skin lesions, suggesting a common inflammatory mechanism.
Gene mutations of proteins involved in the skin's ability to function as a barrier have been identified as markers of susceptibility for the development of psoriasis.
Deoxyribonucleic acid (DNA) released from dying cells acts as an inflammatory stimulus in psoriasis and stimulates the receptors on certain dendritic cells, which in turn produce the cytokine interferon-α. In response to these chemical messages from dendritic cells and T cells, keratinocytes also secrete cytokines such as interleukin-1, interleukin-6, and tumor necrosis factor-α, which signal downstream inflammatory cells to arrive and stimulate additional inflammation.
Dendritic cells bridge the innate immune system and adaptive immune system. They are increased in psoriatic lesions and induce the proliferation of T cells and type 1 helper T cells (Th1). Targeted immunotherapy, as well as psoralen and ultraviolet A (PUVA) therapy, can reduce the number of dendritic cells and favors a TH2 cell cytokine secretion pattern over a Th1/Th17 cell cytokine profile. Psoriatic T cells move from the dermis into the epidermis and secrete interferon-γ and interleukin-17. Interleukin-23 is known to induce the production of interleukin-17 and interleukin-22. Interleukin-22 works in combination with interleukin-17 to induce keratinocytes to secrete neutrophil-attracting cytokines.
Diagnosis
A diagnosis of psoriasis is usually based on the appearance of the skin. Skin characteristics typical for psoriasis are scaly, erythematous plaques, papules, or patches of skin that may be painful and itch. No special blood tests or diagnostic procedures are usually required to make the diagnosis.
The differential diagnosis of psoriasis includes dermatological conditions similar in appearance such as discoid eczema, seborrheic eczema, pityriasis rosea (may be confused with guttate psoriasis), nail fungus (may be confused with nail psoriasis) or cutaneous T cell lymphoma (50% of individuals with this cancer are initially misdiagnosed with psoriasis). Dermatologic manifestations of systemic illnesses such as the rash of secondary syphilis may also be confused with psoriasis.
If the clinical diagnosis is uncertain, a skin biopsy or scraping may be performed to rule out other disorders and to confirm the diagnosis. Skin from a biopsy shows clubbed epidermal projections that interdigitate with dermis on microscopy. Epidermal thickening is another characteristic histologic finding of psoriasis lesions. The stratum granulosum layer of the epidermis is often missing or significantly decreased in psoriatic lesions; the skin cells from the most superficial layer of skin are also abnormal as they never fully mature. Unlike their mature counterparts, these superficial cells keep their nuclei. Inflammatory infiltrates can typically be seen on microscopy when examining skin tissue or joint tissue affected by psoriasis. Epidermal skin tissue affected by psoriatic inflammation often has many CD8+ T cells, while a predominance of CD4+ T cells makes up the inflammatory infiltrates of the dermal layer of skin and joints.
Classification
Morphological
Psoriasis is classified as a papulosquamous disorder and is most commonly subdivided into different categories based on histological characteristics. Variants include plaque, pustular, guttate, and flexural psoriasis. Each form has a dedicated ICD-10 code. Psoriasis can also be classified into nonpustular and pustular types.
Pathogenetic
Another classification scheme considers genetic and demographic factors. Type 1 has a positive family history, starts before the age of 40, and is associated with the human leukocyte antigen, HLA-Cw6. Conversely, type 2 does not show a family history, presents after age 40, and is not associated with HLA-Cw6. Type 1 accounts for about 75% of persons with psoriasis.
The classification of psoriasis as an autoimmune disease has sparked considerable debate. Researchers have proposed differing descriptions of psoriasis and psoriatic arthritis; some authors have classified them as autoimmune diseases while others have classified them as distinct from autoimmune diseases and referred to them as immune-mediated inflammatory diseases.
Severity
No consensus exists about how to classify the severity of psoriasis. Mild psoriasis has been defined as a percentage of body surface area (BSA)≤10, a Psoriasis Area and Severity Index (PASI) score ≤10, and a Dermatology Life Quality Index (DLQI) score ≤10. Moderate to severe psoriasis was defined by the same group as BSA >10 or PASI score >10 and a DLQI score >10.
The DLQI is a 10-question tool used to measure the impact of several dermatologic diseases on daily functioning. The DLQI score ranges from 0 (minimal impairment) to 30 (maximal impairment) and is calculated with each answer being assigned 0–3 points with higher scores indicating greater social or occupational impairment.
The PASI is the most widely used measurement tool for psoriasis. It assesses the severity of lesions and the area affected and combines these two factors into a single score from 0 (no disease) to 72 (maximal disease). Nevertheless, the PASI can be too unwieldy to use outside of research settings, which has led to attempts to simplify the index for clinical use.
Co-morbidities
Psoriasis is not just a skin disease. The symptoms of psoriasis can sometimes go beyond the skin and can have a negative impact on the quality of life of the affected individuals. Additionally, the co-morbidities increase the treatment and financial burden of psoriasis and should be considered when managing this condition.
Cardiovascular complications
There is 2.2 times increased risk of cardiovascular complications in people with psoriasis. Also, people with psoriasis are more susceptible to myocardial infarction (heart attack) and stroke. It has been speculated that there is systemic inflammation in psoriasis, which drives “psoriatic march” and can cause other inflammatory complications including cardiovascular complications. A study used fluorodeoxyglucose F-18 positron emission tomography-computed tomography (FDG PET/CT) to measure aortic vascular inflammation in psoriasis patients, and found increased coronary artery disease indices, including total plaque burden, luminal stenosis, and high-risk plaques in people with psoriasis. Similarly, it was found that there is an 11% reduction in aortic vascular inflammation when there is a 75% reduction in the PASI score.
Depression
Depression or depressive symptoms are present in 28–55% of people with psoriasis. People with psoriasis are often stigmatized due to visible disfigurement of the skin. Social stigmatization is a risk factor for depression, however, other immune system factors may also be related to this observed increased incidence of depression in people with psoriasis. There is some evidence that increased inflammatory signals in the body could also contribute to depression in people with chronic inflammatory diseases, including psoriasis.
Type 2 diabetes
People with psoriasis are at increased risk of developing type 2 diabetes (~1.5 odds ratio). A genome-wide genetic study found that psoriasis and type 2 diabetes share four loci, namely, ACTR2, ERLIN1, TRMT112, and BECN1, which are connected via inflammatory NF-κB pathway.
Management
While no cure is available for psoriasis, many treatment options exist. Topical agents are typically used for mild disease, phototherapy for moderate disease, and systemic agents for severe disease. There is no evidence to support the effectiveness of conventional topical and systemic drugs, biological therapy, or phototherapy for acute guttate psoriasis or an acute guttate flare of chronic psoriasis.
Topical agents
Topical corticosteroid preparations are the most effective agents when used continuously for eight weeks; retinoids and coal tar were found to be of limited benefit and may be no better than placebo. Very potent topical corticosteroids may be helpful in some cases, however, it is suggested to only use them for four weeks at a time and only if other less potent topical treatment options are not working.
Vitamin D analogues (such as paricalcitol, calcipotriol, tacalcitol, and calcitriol) are superior to placebo. Combination therapy with vitamin D and a corticosteroid is superior to either treatment alone and vitamin D is superior to coal tar for chronic plaque psoriasis.
For psoriasis of the scalp, a 2016 review found dual therapy (vitamin D analogs and topical corticosteroids) or corticosteroid monotherapy to be more effective and safer than topical vitamin D analogs alone. Due to their similar safety profiles and minimal benefit of dual therapy over monotherapy, corticosteroid monotherapy appears to be an acceptable treatment for short-term treatment.
Moisturizers and emollients such as mineral oil, petroleum jelly, and decubal (an oil-in-water emollient) were found to increase the clearance of psoriatic plaques. Some emollients are even more effective at clearing psoriatic plaques when combined with phototherapy. Certain emollients, though, have no impact on psoriasis plaque clearance or may even decrease the clearance achieved with phototherapy, e.g. the emollient salicylic acid is structurally similar to para-aminobenzoic acid, commonly found in sunscreen, and is known to interfere with phototherapy in psoriasis. Coconut oil, when used as an emollient in psoriasis, has been found to decrease plaque clearance with phototherapy. Medicated creams and ointments applied directly to psoriatic plaques can help reduce inflammation, remove built-up scale, reduce skin turnover, and clear affected skin of plaques. Ointment and creams containing coal tar, dithranol, corticosteroids (i.e. desoximetasone), fluocinonide, vitamin D3 analogues (for example, calcipotriol), and retinoids are routinely used. (The use of the finger tip unit may be helpful in guiding how much topical treatment to use.)
Vitamin D analogs may be useful with steroids; steroids alone have a higher rate of side effects. Vitamin D analogs may allow lower doses of steroids to be used.
Another topical therapy used to treat psoriasis is a form of balneotherapy, which involves daily baths in saltwater, such as the Dead Sea, combined with sun exposure. This is usually done for four weeks in which exposure time is gradually increased. The primary benefit is attributed to sun exposure and specifically UVB light. This is cost-effective and it has been propagated as an effective way to treat psoriasis without medication. Decreases of PASI scores greater than 75% and remission for several months have commonly been observed. Side effects may be mild such as itchiness, folliculitis, sunburn, poikiloderma, and a theoretical risk of nonmelanoma cancer or melanoma has been suggested. Some studies indicate no increased risk of melanoma in the long term. Data are inconclusive concerning nonmelanoma skin cancer risk, but support the idea that the therapy is associated with an increased risk of benign forms of sun-induced skin damage such as, but not limited to, actinic elastosis or liver spots. Dead Sea balneotherapy is also effective for psoriatic arthritis. Tentative evidence indicates that balneophototherapy, a combination of salt bathes and exposure to ultraviolet B-light (UVB), in chronic plaque psoriasis is better than UVB alone. Glycerin is also an effective treatment for Psoriasis.
UV phototherapy
Phototherapy in the form of sunlight has long been used for psoriasis. UVB wavelengths of 311–313 nanometers are most common. These lamps have been developed for this treatment. The exposure time should be controlled to avoid overexposure and burning of the skin. The UVB lamps should have a timer that turns off the lamp when the time ends. The dose is increased in every treatment to let the skin get used to the light. Increased rates of cancer from treatment appear to be small. Narrowband UVB therapy has been demonstrated to have similar efficacy to psoralen and ultraviolet A phototherapy (PUVA). A 2013 meta-analysis found no difference in efficacy between NB-UVB and PUVA in the treatment of psoriasis, but NB-UVB is usually more convenient.
One of the problems with clinical phototherapy is the difficulty many people have gaining access to a facility. Indoor tanning resources are almost ubiquitous today and could be considered as a means for people to get UV exposure when dermatologist-provided phototherapy is not available. Indoor tanning is already used by many people as a treatment for psoriasis; one indoor facility reported that 50% of its clients were using the center for psoriasis treatment; another reported 36% were doing the same thing. However, a concern with the use of commercial tanning is that tanning beds that primarily emit UVA might not effectively treat psoriasis. One study found that plaque psoriasis is responsive to erythemogenic doses of either UVA or UVB, as exposure to either can cause dissipation of psoriatic plaques. It does require more energy to reach erythemogenic dosing with UVA.
UV light therapies all have risks; tanning beds are no exception, being listed by the World Health Organization as carcinogens. Exposure to UV light is known to increase the risks of melanoma and squamous cell and basal cell carcinomas; younger people with psoriasis, particularly those under age 35, are at increased risk from melanoma from UV light treatment. A review of studies recommends that people who are susceptible to skin cancers exercise caution when using UV light therapy as a treatment.
A major mechanism of NB-UVB is the induction of DNA damage in the form of pyrimidine dimers. This type of phototherapy is useful in the treatment of psoriasis because the formation of these dimers interferes with the cell cycle and stops it. The interruption of the cell cycle induced by NB-UVB opposes the characteristic rapid division of skin cells seen in psoriasis. The activity of many types of immune cells found in the skin is also effectively suppressed by NB-UVB phototherapy treatments. The most common short-term side effect of this form of phototherapy is redness of the skin; less common side effects of NB-UVB phototherapy are itching and blistering of the treated skin, irritation of the eyes in the form of conjunctival inflammation or inflammation of the cornea, or cold sores due to reactivation of the herpes simplex virus in the skin surrounding the lips. Eye protection is usually given during phototherapy treatments.
PUVA combines the oral or topical administration of psoralen with exposure to ultraviolet A (UVA) light. The mechanism of action of PUVA is unknown but probably involves activation of psoralen by UVA light, which inhibits the abnormally rapid production of the cells in psoriatic skin. There are multiple mechanisms of action associated with PUVA, including effects on the skin's immune system. PUVA is associated with nausea, headache, fatigue, burning, and itching. Long-term treatment is associated with squamous cell carcinoma (but not with melanoma). A combination therapy for moderate to severe psoriasis using PUVA plus acitretin resulted in benefit, but acitretin use has been associated with birth defects and liver damage.
Systemic agents
Psoriasis resistant to topical treatment and phototherapy may be treated with systemic therapies including medications by mouth or injectable treatments. People undergoing systemic treatment must have regular blood and liver function tests to check for medication toxicities. Pregnancy must be avoided for most of these treatments. The majority of people experience a recurrence of psoriasis after systemic treatment is discontinued.
Non-biologic systemic treatments frequently used for psoriasis include methotrexate, ciclosporin, hydroxycarbamide, fumarates such as dimethyl fumarate, and retinoids. Methotrexate and ciclosporin are medications that suppress the immune system; retinoids are synthetic forms of vitamin A. These agents are also regarded as first-line treatments for psoriatic erythroderma. Oral corticosteroids should not be used as they can severely flare psoriasis upon their discontinuation.
Biologics are manufactured proteins that interrupt the immune process involved in psoriasis. Unlike generalized immunosuppressive medical therapies such as methotrexate, biologics target specific aspects of the immune system contributing to psoriasis. These medications are generally well-tolerated, and limited long-term outcome data have demonstrated biologics to be safe for long-term use in moderate to severe plaque psoriasis. However, due to their immunosuppressive actions, biologics have been associated with a small increase in the risk for infection.
Guidelines regard biologics as a third-line treatment for plaque psoriasis following inadequate response to topical treatment, phototherapy, and non-biologic systemic treatments. The safety of biologics during pregnancy has not been assessed. European guidelines recommend avoiding biologics if a pregnancy is planned; anti-TNF therapies such as infliximab are not recommended for use in chronic carriers of the hepatitis B virus or individuals infected with HIV.
Several monoclonal antibodies target cytokines, the molecules that cells use to send inflammatory signals to each other. TNF-α is one of the main executor inflammatory cytokines. Four monoclonal antibodies (MAbs) (infliximab, adalimumab, golimumab, and certolizumab pegol) and one recombinant TNF-α decoy receptor, etanercept, have been developed to inhibit TNF-α signaling. Additional monoclonal antibodies, such as ixekizumab, have been developed against pro-inflammatory cytokines and inhibit the inflammatory pathway at a different point than the anti-TNF-α antibodies. IL-12 and IL-23 share a common domain, p40, which is the target of the FDA-approved ustekinumab. In 2017 the US FDA approved guselkumab for plaque psoriasis. There have been few studies of the efficacy of anti-TNF medications for psoriasis in children. One randomized control study suggested that 12 weeks of etanercept treatment reduced the extent of psoriasis in children with no lasting adverse effects.
Two medications that target T cells are efalizumab and alefacept. Efalizumab is a monoclonal antibody that specifically targets the CD11a subunit of LFA-1. It also blocks the adhesion molecules on the endothelial cells that line blood vessels, which attract T cells. Efalizumab was voluntarily withdrawn from the European market in February 2009, and from the U.S. market in June 2009, by the manufacturer due to the medication's association with cases of progressive multifocal leukoencephalopathy. Alefacept also blocks the molecules that dendritic cells use to communicate with T cells and even causes natural killer cells to kill T cells as a way of controlling inflammation. Apremilast may also be used.
Individuals with psoriasis may develop neutralizing antibodies against monoclonal antibodies. Neutralization occurs when an antidrug antibody prevents a monoclonal antibody such as infliximab from binding antigen in a laboratory test. Specifically, neutralization occurs when the anti-drug antibody binds to infliximab's antigen binding site instead of TNF-α. When infliximab no longer binds tumor necrosis factor alpha, it no longer decreases inflammation, and psoriasis may worsen. Neutralizing antibodies have not been reported against etanercept, a biologic medication that is a fusion protein composed of two TNF-α receptors. The lack of neutralizing antibodies against etanercept is probably secondary to the innate presence of the TNF-α receptor, and the development of immune tolerance.
There is strong evidence to indicate that infliximab, bimekizumab, ixekizumab, and risankizumab are the most effective biologics for treating moderate to severe cases of psoriasis. There is also some evidence to support use of secukinumab, brodalumab, guselkumab, certolizumab, and ustekinumab. In general, anti-IL17, anti-IL12/23, anti-IL23, and anti-TNF alpha biologics were found to be more effective than traditional systemic treatments. The immunologic pathways of psoriasis involve Th9, Th17, Th1 lymphocytes, and IL-22. The aforementioned biologic agents hinder different aspects of these pathways.
Another set of treatments for moderate to severe psoriasis are fumaric acid esters (FAE), which may be similar in effectiveness to methotrexate.
Apremilast (Otezla, Celgene) is an oral small-molecule inhibitor of the enzyme phosphodiesterase 4, which plays an important role in chronic inflammation associated with psoriasis.
It has been theorized that antistreptococcal medications may improve guttate and chronic plaque psoriasis; however, limited studies do not show that antibiotics are effective.
Surgery
Limited evidence suggests removal of the tonsils may benefit people with chronic plaque psoriasis, guttate psoriasis, and palmoplantar pustulosis.
Diet
Uncontrolled studies have suggested that individuals with psoriasis or psoriatic arthritis may benefit from a diet supplemented with fish oil rich in eicosapentaenoic acid (EPA) and docosahexaenoic acid (DHA). A low-calorie diet appears to reduce the severity of psoriasis. Diet recommendations include consumption of cold water fish (preferably wild fish, not farmed) such as salmon, herring, and mackerel; extra virgin olive oil; legumes; vegetables; fruits; and whole grains; and avoid consumption of alcohol, red meat, and dairy products (due to their saturated fat). The effect of caffeine consumption (including from coffee, black tea, mate, and dark chocolate) remains to be determined.
Many patients report improvements after consuming less tobacco, caffeine, sugar, nightshades (tomatoes, eggplant, peppers, paprika and white potatoes) and taking probiotics and oral Vitamin D.
There is a higher rate of celiac disease among people with psoriasis. When adopting a gluten-free diet, disease severity generally decreases in people with celiac disease and those with anti-gliadin antibodies.
Prognosis
Most people with psoriasis experience nothing more than mild skin lesions that can be treated effectively with topical therapies. Depending on the severity and location of outbreaks, people may experience significant physical discomfort and some disability, affecting the person's quality of life. Itching and pain can interfere with basic functions, such as self-care and sleep. Participation in sporting activities, certain occupations, and caring for family members can become difficult activities for those with plaques located on their hands and feet. Plaques on the scalp can be particularly embarrassing, as flaky plaque in the hair can be mistaken for dandruff.
thumb|100px|Filipina with psoriasis
Individuals with psoriasis may feel self-conscious about their appearance and have a poor self-image that stems from fear of public rejection and psychosexual concerns. Psoriasis has been associated with low self-esteem and depression is more common among those with the condition. People with psoriasis often feel prejudiced against due to the commonly held incorrect belief that psoriasis is contagious. Psychological distress can lead to significant depression and social isolation; a high rate of thoughts about suicide has been associated with psoriasis. Many tools exist to measure the quality of life of people with psoriasis and other dermatological disorders. Clinical research has indicated individuals often experience a diminished quality of life. Children with psoriasis may encounter bullying.
Several conditions are associated with psoriasis including obesity, cardiovascular, and metabolic disturbances. These occur more frequently in older people. Nearly half of individuals with psoriasis over the age of 65 have at least three comorbidities (concurrent conditions), and two-thirds have at least two comorbidities.
Cardiovascular disease
Psoriasis has been associated with obesity and several other cardiovascular and metabolic disturbances. The number of new cases per year of diabetes is 27% higher in people affected by psoriasis than in those without the condition. Severe psoriasis may be even more strongly associated with the development of diabetes than mild psoriasis. Younger people with psoriasis may also be at increased risk for developing diabetes. Individuals with psoriasis or psoriatic arthritis have a slightly higher risk of heart disease and heart attacks when compared to the general population. Cardiovascular disease risk appeared to be correlated with the severity of psoriasis and its duration. There is no strong evidence to suggest that psoriasis is associated with an increased risk of death from cardiovascular events. Methotrexate may provide a degree of protection for the heart.
The odds of having hypertension are 1.58 times higher in people with psoriasis than those without the condition; these odds are even higher with severe cases of psoriasis. A similar association was noted in people who have psoriatic arthritis—the odds of having hypertension were found to be 2.07 times greater when compared to odds of the general population. The link between psoriasis and hypertension is not currently understood. Mechanisms hypothesized to be involved in this relationship include the following: dysregulation of the renin–angiotensin system, elevated levels of endothelin 1 in the blood, and increased oxidative stress. The number of new cases of the heart rhythm abnormality atrial fibrillation is 1.31 times higher in people with mild psoriasis and 1.63 times higher in people with severe psoriasis. There may be a slightly increased risk of stroke associated with psoriasis, especially in severe cases. Treating high levels of cholesterol with statins has been associated with decreased psoriasis severity, as measured by PASI score, and has also been associated with improvements in other cardiovascular disease risk factors such as markers of inflammation. These cardioprotective effects are attributed to ability of statins to improve blood lipid profile and because of their anti-inflammatory effects. Statin use in those with psoriasis and hyperlipidemia was associated with decreased levels of high-sensitivity C-reactive protein and TNFα as well as decreased activity of the immune protein LFA-1. Compared to individuals without psoriasis, those affected by psoriasis are more likely to satisfy the criteria for metabolic syndrome.
Other diseases
The rates of Crohn's disease and ulcerative colitis are increased when compared with the general population, by a factor of 3.8 and 7.5 respectively. People with psoriasis also have a higher risk of celiac disease. Few studies have evaluated the association of multiple sclerosis with psoriasis, and the relationship has been questioned. Psoriasis has been associated with a 16% increase in overall relative risk for non-skin cancer, thought to be attributed to systemic therapy, particularly methotrexate. People treated with long-term systemic therapy for psoriasis have a 52% increased risk cancers of the lung and bronchus, a 205% increase in the risk of developing cancers of the upper gastrointestinal tract, a 31% increase in the risk of developing cancers of the urinary tract, a 90% increase in the risk of developing liver cancer, and a 46% increase in the risk of developing pancreatic cancer. The risk for development of non-melanoma skin cancers is also increased. Psoriasis increases the risk of developing squamous cell carcinoma of the skin by 431% and increases the risk of basal cell carcinoma by 100%. There is no increased risk of melanoma associated with psoriasis. People with psoriasis have a higher risk of developing cancer.
Epidemiology
Psoriasis is estimated to affect 2–4% of the population of the western world. The rate of psoriasis varies according to age, region and ethnicity; a combination of environmental and genetic factors is thought to be responsible for these differences. Psoriasis is about five times more common in people of European descent than in people of Asian descent, more common in countries farther from the equator, relatively uncommon in African Americans, and extremely uncommon in Native Americans. Psoriasis has been estimated to affect about 6.7million Americans.
Psoriasis can occur at any age, although it is more frequent in adults and commonly appears for the first time between the ages of 15 and 25 years. Approximately one-third of people with psoriasis report being diagnosed before age 20. Psoriasis affects both sexes equally.
People with inflammatory bowel disease such as Crohn's disease or ulcerative colitis are at an increased risk of developing psoriasis.
History
Scholars believe psoriasis to have been included among the various skin conditions called tzaraath (translated as leprosy) in the Hebrew Bible. The person was deemed "impure" (see tumah and taharah) during their affected phase and is ultimately treated by the kohen. However, it is more likely that this confusion arose from the use of the same Greek term for both conditions. The Greeks used the term lepra (λέπρα) for scaly skin conditions. They used the term psora (ψώρα) to describe itchy skin conditions. It became known as Willan's lepra in the late 18th century when English dermatologists Robert Willan and Thomas Bateman differentiated it from other skin diseases. Leprosy, they said, is distinguished by the regular, circular form of patches, while psoriasis is always irregular. Willan identified two categories: leprosa graecorum and psora leprosa.
Psoriasis is thought to have first been described in Ancient Rome by Cornelius Celsus. The British dermatologist Thomas Bateman described a possible link between psoriasis and arthritic symptoms in 1813. Admiral William Halsey missed out on the Battle of Midway because he contracted psoriasis while out at sea in the early months of American participation of World War II. Admiral Chester Nimitz medically ordered Halsey to recover at a hospital in Hawaii.
The history of psoriasis is littered with treatments of dubious effectiveness and high toxicity. In the 18th and 19th centuries, Fowler's solution, which contains a poisonous and carcinogenic arsenic compound, was used by dermatologists as a treatment for psoriasis. Mercury was also used for psoriasis treatment during this time. Sulfur, iodine, and phenol were also commonly used treatments for psoriasis during this era when it was incorrectly believed that psoriasis was an infectious disease. Coal tars were widely used with ultraviolet light irradiation as a topical treatment approach in the early 1900s. During the same time, psoriatic arthritis cases were treated with intravenously administered gold preparations in the same manner as rheumatoid arthritis.
Society and culture
The International Federation of Psoriasis Associations (IFPA) is the global umbrella organization for national and regional psoriasis associations and also gathers the leading experts in psoriasis and psoriatic arthritis research for scientific conferences every three years. The Psoriasis International Network, a program of the Fondation René Touraine, gathers dermatologists, rheumatologists, and other caregivers involved in the management of psoriasis. Non-profit organizations like the National Psoriasis Foundation in the United States, the Psoriasis Association in the United Kingdom, and Psoriasis Australia offer advocacy and education about psoriasis in their respective countries.
Cost
The annual cost of treating psoriasis in the United States is estimated as high as $32.5billion, including $12.2billion in direct costs. Pharmacy costs are the main source of direct expense, with biologic therapy the most prevalent. These costs increase significantly when co-morbid conditions such as heart disease, hypertension, diabetes, lung disease, and psychiatric disorders are factored in. Expenses linked to co-morbidities are estimated at an additional $23,000 per person per year.
Research
The role of insulin resistance in the pathogenesis of psoriasis is under investigation. Preliminary research has suggested that antioxidants such as polyphenols may have beneficial effects on the inflammation characteristic of psoriasis.
Many novel medications being researched during the 2010s target the Th17/IL-23 axis, particularly IL-23p19 inhibitors, as IL-23p19 is present in increased concentrations in psoriasis skin lesions while contributing less to protection against opportunistic infections. Other cytokines such as IL-17 and IL-22 also have been targets for inhibition as they play important roles in the pathogenesis of psoriasis. Another avenue of research has focused on the use of vascular endothelial growth factor inhibitors to treat psoriasis. Oral agents being investigated during the 2010s as alternatives to medications administered by injection include Janus kinase inhibitors, protein kinase C inhibitors, mitogen-activated protein kinase inhibitors, and phosphodiesterase 4 inhibitors, all of which have proven effective in various phase 2 and 3 clinical trials. These agents have potentially severe side-effects due to their immunosuppressive mechanisms.
| Biology and health sciences | Specific diseases | Health |
101970 | https://en.wikipedia.org/wiki/Tinnitus | Tinnitus | Tinnitus is a condition when a person hears a ringing sound or a different variety of sound when no corresponding external sound is present and other people cannot hear it. Nearly everyone experiences faint "normal tinnitus" in a completely quiet room; but this is of concern only if it is bothersome, interferes with normal hearing, or is associated with other problems. The word tinnitus comes from the Latin tinnire, "to ring". In some people, it interferes with concentration, and can be associated with anxiety and depression.
Tinnitus is usually associated with hearing loss and decreased comprehension of speech in noisy environments. It is common, affecting about 10–15% of people. Most tolerate it well, and it is a significant problem in only 1–2% of people. It can trigger a fight-or-flight response, as the brain may perceive it as dangerous and important.
Rather than a disease, tinnitus is a symptom that may result from a variety of underlying causes and may be generated at any level of the auditory system as well as outside that system. The most common causes are hearing damage, noise-induced hearing loss, or age-related hearing loss, known as presbycusis. Other causes include ear infections, disease of the heart or blood vessels, Ménière's disease, brain tumors, acoustic neuromas (tumors on the auditory nerves of the ear), migraines, temporomandibular joint disorders, exposure to certain medications, a previous head injury, and earwax. It can suddenly emerge during a period of emotional stress. It is more common in those with depression.
The diagnosis of tinnitus is usually based on a patient's description of the symptoms they are experiencing. Such a diagnosis is commonly supported by an audiogram, and an otolaryngological and neurological examination. How much tinnitus interferes with a person's life may be quantified with questionnaires. If certain problems are found, medical imaging, such as magnetic resonance imaging (MRI), may be performed. Other tests are suitable when tinnitus occurs with the same rhythm as the heartbeat. Rarely, the sound may be heard by someone other than the patient by using a stethoscope, in which case it is known as "objective tinnitus". Occasionally, spontaneous otoacoustic emissions, sounds produced normally by the inner ear, may result in tinnitus.
Measures to prevent tinnitus include avoiding chronic or extended exposure to loud noise, and limiting exposure to ototoxic drugs and substances. If there is an underlying cause, treating that cause may lead to improvements. Otherwise, typically, tinnitus management involves psychoeducation or counseling, such as talk therapy. Sound generators or hearing aids may help. No medication directly targets tinnitus.
Signs and symptoms
Tinnitus is often described as ringing, but it may also sound like clicking, buzzing, hissing, or roaring. It may be soft or loud, low- or high-pitched, and may seem to come from either one or both ears, or from the head itself. It may be intermittent or continuous. In some individuals, its intensity may be changed by shoulder, neck, head, tongue, jaw, or eye movements.
Course
Due to variations in study designs, data on the course of tinnitus shows few consistent results. Generally, prevalence increases with age in adults, and the ratings of annoyance increase with persistence at follow up.
Psychological effects
Although it is an annoying condition to which most people adapt, persistent tinnitus may cause anxiety and depression in some people. Tinnitus annoyance is more strongly associated with the psychological condition of the person than the loudness or frequency range of the perceived sound. Psychological problems such as depression, anxiety, sleep disturbances, and concentration difficulties are common in those with strongly annoying tinnitus. 45% of people with tinnitus have an anxiety disorder at some time in their lives.
Psychological research has focused on the tinnitus distress reaction to account for differences in tinnitus severity. The research indicates that conditioning at the initial perception of tinnitus linked it with negative emotions, such as fear and anxiety.
Types
Commonly tinnitus is classified into "subjective and objective tinnitus". Tinnitus is usually subjective, meaning that the sounds the person hears are not detectable by means currently available to physicians and hearing technicians. Subjective tinnitus has also been called "tinnitus aurium", "non-auditory", or "non-vibratory" tinnitus. In rare cases, tinnitus can be heard by someone else using a stethoscope. Even more rarely, in some cases it can be measured as a spontaneous otoacoustic emission (SOAE) in the ear canal. This is classified as objective tinnitus, also called "pseudo-tinnitus" or "vibratory" tinnitus.
Subjective tinnitus
Subjective tinnitus is the most frequent type. It can have many causes, but most commonly it results from hearing loss. When it is caused by disorders of the inner ear or auditory nerve, it can be called "otic" (from the Greek word for ear). These otological or neurological disorders include those triggered by infections, drugs, or trauma. A cause is traumatic noise exposure that damages hair cells in the inner ear. Some evidence suggests that long-term exposure to noise pollution from heavy traffic may increase the risk of developing tinnitus.
When there does not seem to be a connection with a disorder of the inner ear or auditory nerve, tinnitus can be called "non-otic". In 30% of cases, tinnitus is influenced by the somatosensory system; for instance, people can increase or decrease their tinnitus by moving their face, head, jaw, or neck. This type is called somatic or craniocervical tinnitus, since it is only head or neck movements that have an effect.
Some tinnitus may be caused by neuroplastic changes in the central auditory pathway. In this theory, the disturbance of sensory input caused by hearing loss results in such changes as a homeostatic response of neurons in the central auditory system, causing tinnitus. When some frequencies of sound are lost to hearing loss, the auditory system compensates by amplifying those frequencies, eventually producing sound sensations at those frequencies constantly even when there is no corresponding external sound.
Hearing loss
The most common cause of tinnitus is hearing loss. Hearing loss may have many different causes, but among those with tinnitus, the major cause is cochlear injury.
In many cases no underlying cause is identified.
Ototoxic drugs also may cause subjective tinnitus, as they may cause hearing loss, or increase the damage done by exposure to loud noise. This damage may occur even at doses not considered ototoxic. More than 260 medications have been reported to cause tinnitus as a side effect.
Tinnitus can also occur from the discontinuation of therapeutic doses of benzodiazepines. It can sometimes be a protracted symptom of benzodiazepine withdrawal and may persist for many months. Medications such as bupropion may also cause tinnitus.
Associated factors
Factors associated with tinnitus include:
Ear problems and hearing loss:
Conductive hearing loss
Acoustic shock
Loud noise or music
Middle ear effusion
Otitis
Otosclerosis
Eustachian tube dysfunction
Sensorineural hearing loss
Excessive or loud noise; e.g. acoustic trauma
Presbycusis (age-associated hearing loss)
Ménière's disease
Endolymphatic hydrops
Superior canal dehiscence
Acoustic neuroma
Mercury or lead poisoning
Ototoxic medications
Neurologic disorders:
Chiari malformation
Multiple sclerosis
Head injury
Giant cell arteritis
Temporomandibular joint dysfunction
Metabolic disorders:
Vitamin B12 deficiency
Iron deficiency anemia
Psychiatric disorders
Depression
Anxiety disorders
Other factors:
Vasculitis
Some psychedelic drugs can produce temporary tinnitus-like symptoms as a side effect:
5-MeO-DET
Diisopropyltryptamine (DiPT)
Benzodiazepine withdrawal
Intracranial hyper or hypotension caused by, for example, encephalitis or a cerebrospinal fluid leak
Objective tinnitus
A specific type of tinnitus, objective tinnitus, is characterized by hearing the sounds of one's own muscle contractions or pulse, typically a result of sounds that have been created by the movement of jaw muscles or sounds related to blood flow in the neck or face. It is sometimes caused by an involuntary twitching of a muscle or a group of muscles (myoclonus) or by a vascular condition. In some cases, tinnitus is generated by muscle spasms around the middle ear.
Spontaneous otoacoustic emissions (SOAEs)—faint high-frequency tones that are produced in the inner ear and can be measured in the ear canal with a sensitive microphone—may also cause tinnitus. About 8% of those with SOAEs and tinnitus have SOAE-linked tinnitus, while the percentage of all cases of tinnitus caused by SOAEs is estimated at 4%.
Pediatric tinnitus
Children may be subject to pulsatile or continuous tinnitus, involving anomalies and variants of the vascular parts affecting the middle/inner ear structures. CT scans may be used to check the integrity of the structures, and MR scans can evaluate nerves and potential masses or malformations. Early diagnosis can prevent long-term impairments to development.
Pulsatile tinnitus
Some people experience a sound that beats in time with their pulse, known as pulsatile tinnitus or vascular tinnitus. Pulsatile tinnitus is usually objective in nature, resulting from altered blood flow or increased blood turbulence near the ear, such as from atherosclerosis or venous hum, but it can also arise as a subjective phenomenon from an increased awareness of blood flow in the ear.
The differential diagnosis for pulsatile tinnitus is wide and includes vascular etiologies, tumors, disorders of the middle ear or inner ear, and other intracranial pathologies. Vascular causes of pulsatile tinnitus include venous causes (e.g., high riding or dehiscent jugular bulb, sigmoid sinus diverticulum), arterial causes (e.g., cervical atherosclerosis, potentially life-threatening conditions such as carotid artery aneurysm or carotid artery dissection), or dural arteriovenous fistula or arteriovenous malformations.
Pulsatile tinnitus may also indicate vasculitis, or more specifically, giant cell arteritis. Pulsatile tinnitus may also be caused by tumors such as paragangliomas (e.g., glomus tympanicum, glomus jugulare) or hemangiomas (e.g., facial nerve or cavernous). Middle ear causes of pulsatile tinnitus include patulous eustachian tube, otosclerosis, or middle ear myoclonus (e.g., stapedial or tensor tympani myoclonus). The most common inner ear cause of pulsatile tinnitus is superior semicircular canal dehiscence. Pulsatile tinnitus may also indicate idiopathic intracranial hypertension. Pulsatile tinnitus can be a symptom of intracranial vascular abnormalities and should be evaluated for irregular noises of blood flow (bruits).
Pathophysiology
Tinnitus may be caused by increased neural activity in the auditory brainstem, where the brain processes sounds, causing some auditory nerve cells to become overexcited. The basis of this theory is that many with tinnitus also have hearing loss.
Three reviews in 2016 emphasized the large range and possible combinations of pathologies involved in tinnitus, which result in a great variety of symptoms and specifically adapted therapies.
Diagnosis
The diagnostic approach is based on a history of the condition and an examination of the head, neck, and neurological system. Typically an audiogram is done, and occasionally medical imaging or electronystagmography. Treatable conditions may include middle ear infection, acoustic neuroma,
concussion, and otosclerosis.
Evaluation of tinnitus can include a hearing test (audiogram), measurement of acoustic parameters of the tinnitus like pitch and loudness, and psychological assessment of comorbid conditions like depression, anxiety, and stress that are associated with severity of the tinnitus.
One definition of tinnitus, in contrast to normal ear noise experience, is that tinnitus lasts five minutes at least twice a week. However, people with tinnitus often experience the noise more frequently than this. Tinnitus can be present constantly or intermittently. Some people with constant tinnitus might not be aware of it all the time, but only, for example, during the night when there is less environmental noise to mask it. Chronic tinnitus can be defined as tinnitus with a duration of six months or more.
Audiology
Since most people with tinnitus also have hearing loss, a pure tone hearing test resulting in an audiogram may help diagnose a cause. An audiogram may also facilitate fitting of a hearing aid in those cases where hearing loss is significant. The pitch of tinnitus is often in the range of the hearing loss.
Psychoacoustics
Acoustic qualification of tinnitus includes measurement of several acoustic parameters like frequency in cases of monotone tinnitus or , loudness in dB above hearing threshold at the indicated frequency, , and minimum masking level. In most cases, tinnitus pitch or frequency range is between 5 kHz and 10 kHz, and loudness between 5 and 15 dB above the hearing threshold.
Another relevant parameter of tinnitus is residual inhibition: the temporary suppression or disappearance of tinnitus following a period of masking. The degree of residual inhibition may indicate how effective tinnitus maskers would be as treatment.
An assessment of hyperacusis, a frequent accompaniment of tinnitus, may also be made. Hyperacusis is related to negative reactions to sound and can take many forms. One parameter that can be measured is Loudness Discomfort Level (LDL) in dB, the subjective level of acute discomfort at specified frequencies over the frequency range of hearing. This defines a dynamic range between the hearing threshold at that frequency and the loudness discomfort level. A compressed dynamic range over a particular frequency range can be associated with hyperacusis. Normal hearing threshold is generally defined as 0–20 decibels (dB). Normal loudness discomfort levels are 85–90+ dB, with some authorities citing 100 dB. A dynamic range of 55 dB or less is indicative of hyperacusis.
Severity
Tinnitus is often rated on a scale from "slight" to "severe" according to the effects it has, such as interference with sleep, quiet activities, and normal daily activities.
Assessment of psychological processes related to tinnitus involves measurement of tinnitus severity and distress, as measured subjectively by validated self-report tinnitus questionnaires. Such questionnaires measure the degree of psychological distress and handicap associated with tinnitus, including effects on hearing, lifestyle, health, and emotional functioning. A broader assessment of general functioning, such as levels of anxiety, depression, stress, life stressors, and sleep difficulties, is also important in the assessment of tinnitus due to higher risk of negative well-being across these areas, which may be affected by or exacerbate the tinnitus symptoms.
Current assessment measures aim to identify levels of distress and interference, coping responses, and perceptions of tinnitus to inform treatment and monitor progress. However, wide variability, inconsistencies, and lack of consensus regarding assessment methodology are evidenced in the literature, limiting comparison of treatment effectiveness. Developed to guide diagnosis or classify severity, most tinnitus questionnaires have been shown to be treatment-sensitive outcome measures.
Pulsatile tinnitus
If examination reveals a bruit (sound due to turbulent blood flow), imaging studies such as transcranial doppler (TCD) or magnetic resonance angiography (MRA) should be performed.
Differential diagnosis
Other potential sources of the sounds normally associated with tinnitus should be ruled out. For instance, two recognized sources of high-pitched sounds might be electromagnetic fields common in modern wiring and various sound signal transmissions. A common and often misdiagnosed condition that mimics tinnitus is radio frequency (RF) hearing, in which subjects hear objectively audible high-pitched transmission frequencies that sound similar to tinnitus.
Prevention
Prolonged exposure to loud sound or noise levels can lead to tinnitus. Custom made ear plugs or other measures can help with prevention. Employers may use hearing loss prevention programs to help educate and prevent dangerous levels of exposure to noise. Government organizations set regulations to ensure employees, if following the protocol, should have minimal risk to permanent damage to their hearing.
Certain groups are advised to wear ear plugs to avoid the risk of tinnitus, such as that caused by overexposure to loud noises like wind noise for motorcycle riders. This includes military personnel, musicians, DJs, agricultural workers, and construction workers as people in those occupations are at a greater risk compared to the general population.
Several medicines have ototoxic effects, which can have a cumulative effect that increases the damage done by noise. If ototoxic medications must be administered, close attention by the physician to prescription details, such as dose and dosage interval, can reduce the damage done.
Management
If a specific underlying cause is determined, treating it may lead to improvements. Otherwise, the primary treatment for tinnitus is talk therapy, sound therapy, or hearing aids. There are no effective drugs that treat tinnitus.
Psychological
The best-supported treatment for tinnitus is cognitive behavioral therapy (CBT). It decreases the stress those with tinnitus feel. This appears to be independent of any effect on depression or anxiety. Acceptance and commitment therapy (ACT) also shows promise in the treatment of tinnitus. Relaxation techniques may also help. A clinical protocol called Progressive Tinnitus Management has been developed by the United States Department of Veterans Affairs.
Sound-based interventions
The application of sound therapy by either hearing aids or tinnitus maskers may help the brain ignore the specific tinnitus frequency. Although these methods are poorly supported by evidence, there are no negative effects. There are several approaches for tinnitus sound therapy. The first is sound modification to compensate for the individual's hearing loss. The second is tailored music therapy, notched at the tinnitus frequency, which may affect lateral inhibition of the notched neural region, suppressing tinnitus.
There is some tentative evidence supporting tinnitus retraining therapy, which aims to reduce tinnitus-related neuronal activity. An alternative tinnitus treatment uses mobile applications that include various methods including masking, sound therapy, and relaxation exercises. Such applications can work as a separate device or as a hearing aid control system.
Neuromonics is another sound-based intervention. Its protocol follows the principle of systematic desensitization and involves a structured rehabilitation program lasting 12 months. Neuromonics therapy employs customized sound signals delivered through a device worn by the patient, which aims to target the specific frequency range associated with their tinnitus perception.
Physical therapy
Physical therapy for tinnitus focuses on relaxing jaw and neck muscles that may contribute to symptoms. Muscle tension, particularly in the jaw muscles like the masseter and medial pterygoid, can radiate to the ears, leading to somatic tinnitus. Specialized physical therapists use neuromuscular techniques to alleviate tension in these areas, which may reduce tinnitus intensity and associated pain in connected areas, such as the jaw, teeth, and ears.
Medications
there were no medications effective for idiopathic tinnitus. There is not enough evidence to determine if antidepressants or acamprosate are useful. There are conflicting studies regarding the effectiveness of benzodiazepines for tinnitus. Usefulness of melatonin, as of 2015, is unclear. It is unclear if anticonvulsants are useful for treating tinnitus. Steroid injections into the middle ear also do not seem to be effective. There is no evidence to suggest that the use of betahistine to treat tinnitus is effective.
Botulinum toxin injection has succeeded in some of the rare cases of objective tinnitus from a palatal tremor.
Caroverine is used in a few countries to treat tinnitus. The evidence for its usefulness is very weak.
Neuromodulation
In 2020, information about clinical trials indicated that bimodal neuromodulation may reduce the symptoms of tinnitus. It is a noninvasive technique that involves applying an electrical stimulus to the tongue while also administering sounds. Equipment associated with the treatments is available through physicians. Studies with it and similar devices continue in several research centers.
In March 2023, the US Food and Drug Administration (FDA) approved Neuromod's Lenire device as a treatment option for tinnitus. In June 2024, the US Department of Veterans Affairs (VA) announced it would begin offering the treatment to veterans with tinnitus, making it the first bimodal neuromodulation device to be awarded a Federal Supply Schedule (FSS) contract from the US Government.
Some evidence supports neuromodulation techniques such as transcranial magnetic stimulation, transcranial direct current stimulation, and neurofeedback.
Alternative medicine
Ginkgo biloba does not appear to be effective. The American Academy of Otolaryngology recommends against taking melatonin or zinc supplements to relieve symptoms of tinnitus, and reported that evidence for the efficacy of many dietary supplements (such as lipoflavonoids, garlic, traditional Chinese/Korean herbal medicine, honeybee larvae, and various other vitamins and minerals, as well as homeopathic preparations) did not exist. A 2016 Cochrane Review also concluded that evidence was not sufficient to support taking zinc supplements to reduce symptoms associated with tinnitus.
Prognosis
While there is no cure, most people with tinnitus get used to it over time; for a minority, it remains a significant problem.
Epidemiology
Adults
Tinnitus affects 1015% of people. About a third of North Americans over 55 experience it. It affects one third of adults at some time in their lives, whereas 10–15% are disturbed enough to seek medical evaluation.
70 million people in Europe are estimated to have tinnitus.
Children
Tinnitus is commonly thought of as a symptom of adulthood, and is often overlooked in children. Children with hearing loss have a high incidence of pediatric tinnitus, even though they do not express the condition or its effect on their lives. Children do not generally report tinnitus spontaneously and their complaints may not be taken seriously. Among those who do complain, there is an increased likelihood of associated otological or neurological pathology such as migraine, juvenile Meniere's disease, or chronic suppurative otitis media. Its reported prevalence varies from 12 to 36% in children with normal hearing thresholds, and up to 66% in children with a hearing loss. Approximately 3–10% of children have been reported to be troubled by tinnitus.
| Biology and health sciences | Symptoms and signs | Health |
102024 | https://en.wikipedia.org/wiki/Wetland | Wetland | A wetland is a distinct semi-aquatic ecosystem whose groundcovers are flooded or saturated in water, either permanently, for years or decades, or only seasonally. Flooding results in oxygen-poor (anoxic) processes taking place, especially in the soils. Wetlands form a transitional zone between waterbodies and dry lands, and are different from other terrestrial or aquatic ecosystems due to their vegetation's roots having adapted to oxygen-poor waterlogged soils. They are considered among the most biologically diverse of all ecosystems, serving as habitats to a wide range of aquatic and semi-aquatic plants and animals, with often improved water quality due to plant removal of excess nutrients such as nitrates and phosphorus.
Wetlands exist on every continent, except Antarctica. The water in wetlands is either freshwater, brackish or saltwater. The main types of wetland are defined based on the dominant plants and the source of the water. For example, marshes are wetlands dominated by emergent herbaceous vegetation such as reeds, cattails and sedges. Swamps are dominated by woody vegetation such as trees and shrubs (although reed swamps in Europe are dominated by reeds, not trees). Mangrove forest are wetlands with mangroves, halophytic woody plants that have evolved to tolerate salty water.
Examples of wetlands classified by the sources of water include tidal wetlands, where the water source is ocean tides; estuaries, water source is mixed tidal and river waters; floodplains, water source is excess water from overflowed rivers or lakes; and bogs and vernal ponds, water source is rainfall or meltwater. The world's largest wetlands include the Amazon River basin, the West Siberian Plain, the Pantanal in South America, and the Sundarbans in the Ganges-Brahmaputra delta.
Wetlands contribute many ecosystem services that benefit people. These include for example water purification, stabilization of shorelines, storm protection and flood control. In addition, wetlands also process and condense carbon (in processes called carbon fixation and sequestration), and other nutrients and water pollutants. Wetlands can act as a sink or a source of carbon, depending on the specific wetland. If they function as a carbon sink, they can help with climate change mitigation. However, wetlands can also be a significant source of methane emissions due to anaerobic decomposition of soaked detritus, and some are also emitters of nitrous oxide.
Humans are disturbing and damaging wetlands in many ways, including oil and gas extraction, building infrastructure, overgrazing of livestock, overfishing, alteration of wetlands including dredging and draining, nutrient pollution, and water pollution. Wetlands are more threatened by environmental degradation than any other ecosystem on Earth, according to the Millennium Ecosystem Assessment from 2005. Methods exist for assessing wetland ecological health. These methods have contributed to wetland conservation by raising public awareness of the functions that wetlands can provide. Since 1971, work under an international treaty seeks to identify and protect "wetlands of international importance."
Definitions and terminology
Technical definitions
A simplified definition of wetland is "an area of land that is usually saturated with water". More precisely, wetlands are areas where "water covers the soil, or is present either at or near the surface of the soil all year or for varying periods of time during the year, including during the growing season". A patch of land that develops pools of water after a rain storm would not necessarily be considered a "wetland", even though the land is wet. Wetlands have unique characteristics: they are generally distinguished from other water bodies or landforms based on their water level and on the types of plants that live within them. Specifically, wetlands are characterized as having a water table that stands at or near the land surface for a long enough period each year to support aquatic plants.
A more concise definition is a community composed of hydric soil and hydrophytes.
Wetlands have also been described as ecotones, providing a transition between dry land and water bodies. Wetlands exist "...at the interface between truly terrestrial ecosystems and aquatic systems, making them inherently different from each other, yet highly dependent on both."
In environmental decision-making, there are subsets of definitions that are agreed upon to make regulatory and policy decisions.
Under the Ramsar international wetland conservation treaty, wetlands are defined as follows:
Article 1.1: "...wetlands are areas of marsh, fen, peatland or water, whether natural or artificial, permanent or temporary, with water that is static or flowing, fresh, brackish or salt, including areas of marine water the depth of which at low tide does not exceed six meters."
Article 2.1: "[Wetlands] may incorporate riparian and coastal zones adjacent to the wetlands, and islands or bodies of marine water deeper than six meters at low tide lying within the wetlands."
An ecological definition of a wetland is "an ecosystem that arises when inundation by water produces soils dominated by anaerobic and aerobic processes, which, in turn, forces the biota, particularly rooted plants, to adapt to flooding".
Sometimes a precise legal definition of a wetland is required. The definition used for regulation by the United States government is: 'The term "wetlands" means those areas that are inundated or saturated by surface or ground water at a frequency and duration to support, and that under normal circumstances do support, a prevalence of vegetation typically adapted for life in saturated soil conditions. Wetlands generally included swamps, marshes, bogs, and similar areas.'
For each of these definitions and others, regardless of the purpose, hydrology is emphasized (shallow waters, water-logged soils). The soil characteristics and the plants and animals controlled by the wetland hydrology are often additional components of the definitions.
Types
Wetlands can be tidal (inundated by tides) or non-tidal. The water in wetlands is either freshwater, brackish, saline, or alkaline. There are four main kinds of wetlands – marsh, swamp, bog, and fen (bogs and fens being types of peatlands or mires). Some experts also recognize wet meadows and aquatic ecosystems as additional wetland types. Sub-types include mangrove forests, carrs, pocosins, floodplains, peatlands, vernal pools, sinks, and many others.
The following three groups are used within Australia to classify wetland by type: Marine and coastal zone wetlands, inland wetlands and human-made wetlands. In the US, the best known classifications are the Cowardin classification system and the hydrogeomorphic (HGM) classification system. The Cowardin system includes five main types of wetlands: marine (ocean-associated), estuarine (mixed ocean- and river-associated), riverine (within river channels), lacustrine (lake-associated) and palustrine (inland nontidal habitats).
Peatlands
Peatlands are a unique kind of wetland where lush plant growth and slow decay of dead plants (under anoxic conditions) results in organic peat accumulating; bogs, fens, and mires are different names for peatlands.
Wetland names
Variations of names for wetland systems:
Bayou
Flooded grasslands and savannas
Marsh
Brackish marsh
Freshwater marsh
Mire
Fen
Bog
Riparian zone
Swamp
Freshwater swamp forest
Tidal Freshwater forest
Coniferous swamp
Peat swamp forest
Mangrove swamp
Vernal pool
Some wetlands have localized names unique to a region such as the prairie potholes of North America's northern plain, pocosins, Carolina bays and baygalls of the Southeastern US, mallines of Argentina, Mediterranean seasonal ponds of Europe and California, turloughs of Ireland, billabongs of Australia, among many others.
Locations
By temperature zone
Wetlands are found throughout the world in different climates. Temperatures vary greatly depending on the location of the wetland. Many of the world's wetlands are in the temperate zones, midway between the North or South Poles and the equator. In these zones, summers are warm and winters are cold, but temperatures are not extreme. In subtropical zone wetlands, such as along the Gulf of Mexico, average temperatures might be . Wetlands in the tropics are subjected to much higher temperatures for a large portion of the year. Temperatures for wetlands on the Arabian Peninsula can exceed and these habitats would therefore be subject to rapid evaporation. In northeastern Siberia, which has a polar climate, wetland temperatures can be as low as . Peatlands in arctic and subarctic regions insulate the permafrost, thus delaying or preventing its thawing during summer, as well as inducing its formation.
By precipitation amount
The amount of precipitation a wetland receives varies widely according to its area. Wetlands in Wales, Scotland, and western Ireland typically receive about per year. In some places in Southeast Asia, where heavy rains occur, they can receive up to . In some drier regions, wetlands exist where as little as precipitation occurs each year.
Temporal variation:
Perennial systems
Seasonal systems
Episodic (periodic or intermittent) systems
Ephemeral (short-lived) systems
Surface flow may occur in some segments, with subsurface flow in other segments.
Processes
Wetlands vary widely due to local and regional differences in topography, hydrology, vegetation, and other factors, including human involvement. Other important factors include fertility, natural disturbance, competition, herbivory, burial and salinity. When peat accumulates, bogs and fens arise.
Hydrology
The most important factor producing wetlands is hydrology, or flooding. The duration of flooding or prolonged soil saturation by groundwater determines whether the resulting wetland has aquatic, marsh or swamp vegetation. Other important factors include soil fertility, natural disturbance, competition, herbivory, burial, and salinity. When peat from dead plants accumulates, bogs and fens develop.
Wetland hydrology is associated with the spatial and temporal dispersion, flow, and physio-chemical attributes of surface and ground waters. Sources of hydrological flows into wetlands are predominantly precipitation, surface water (saltwater or freshwater), and groundwater. Water flows out of wetlands by evapotranspiration, surface flows and tides, and subsurface water outflow. Hydrodynamics (the movement of water through and from a wetland) affects hydro-periods (temporal fluctuations in water levels) by controlling the water balance and water storage within a wetland.
Landscape characteristics control wetland hydrology and water chemistry. The O2 and CO2 concentrations of water depend upon temperature, atmospheric pressure and mixing with the air (from winds or water flows). Water chemistry within wetlands is determined by the pH, salinity, nutrients, conductivity, soil composition, hardness, and the sources of water. Water chemistry varies across landscapes and climatic regions. Wetlands are generally minerotrophic (waters contain dissolved materials from soils) with the exception of ombrotrophic bogs that are fed only by water from precipitation.
Because bogs receive most of their water from precipitation and humidity from the atmosphere, their water usually has low mineral ionic composition. In contrast, wetlands fed by groundwater or tides have a higher concentration of dissolved nutrients and minerals.
Fen peatlands receive water both from precipitation and ground water in varying amounts so their water chemistry ranges from acidic with low levels of dissolved minerals to alkaline with high accumulation of calcium and magnesium.
Role of salinity
Salinity has a strong influence on wetland water chemistry, particularly in coastal wetlands and in arid and semiarid regions with large precipitation deficits. Natural salinity is regulated by interactions between ground and surface water, which may be influenced by human activity.
Soil
Carbon is the major nutrient cycled within wetlands. Most nutrients, such as sulfur, phosphorus, carbon, and nitrogen are found within the soil of wetlands. Anaerobic and aerobic respiration in the soil influences the nutrient cycling of carbon, hydrogen, oxygen, and nitrogen, and the solubility of phosphorus thus contributing to the chemical variations in its water. Wetlands with low pH and saline conductivity may reflect the presence of acid sulfates and wetlands with average salinity levels can be heavily influenced by calcium or magnesium. Biogeochemical processes in wetlands are determined by soils with low redox potential.
Biology
The life forms of a wetland system includes its plants (flora) and animals (fauna) and microbes (bacteria, fungi). The most important factor is the wetland's duration of flooding. Other important factors include fertility and salinity of the water or soils. The chemistry of water flowing into wetlands depends on the source of water, the geological material that it flows through and the nutrients discharged from organic matter in the soils and plants at higher elevations. Plants and animals may vary within a wetland seasonally or in response to flood regimes.
Flora
There are four main groups of hydrophytes that are found in wetland systems throughout the world.
Submerged wetland vegetation can grow in saline and fresh-water conditions. Some species have underwater flowers, while others have long stems to allow the flowers to reach the surface. Submerged species provide a food source for native fauna, habitat for invertebrates, and also possess filtration capabilities. Examples include seagrasses and eelgrass.
Floating water plants or floating vegetation are usually small, like those in the Lemnoideae subfamily (duckweeds).
Emergent vegetation like the cattails (Typha spp.), sedges (Carex spp.) and arrow arum (Peltandra virginica) rise above the surface of the water.
When trees and shrubs comprise much of the plant cover in saturated soils, those areas in most cases are called swamps. The upland boundary of swamps is determined partly by water levels. This can be affected by dams Some swamps can be dominated by a single species, such as silver maple swamps around the Great Lakes. Others, like those of the Amazon basin, have large numbers of different tree species. Other examples include cypress (Taxodium) and mangrove swamps.
Fauna
Many species of fish are highly dependent on wetland ecosystems. Seventy-five percent of the United States' commercial fish and shellfish stocks depend solely on estuaries to survive.
Amphibians such as frogs and salamanders need both terrestrial and aquatic habitats in which to reproduce and feed. Because amphibians often inhabit depressional wetlands like prairie potholes and Carolina bays, the connectivity among these isolated wetlands is an important control of regional populations. While tadpoles feed on algae, adult frogs forage on insects. Frogs are sometimes used as an indicator of ecosystem health because their thin skin permits absorption of nutrients and toxins from the surrounding environment resulting in increased extinction rates in unfavorable and polluted environmental conditions.
Reptiles such as snakes, lizards, turtles, alligators and crocodiles are common in wetlands of some regions. In freshwater wetlands of the Southeastern US, alligators are common and a freshwater species of crocodile occurs in South Florida. The Florida Everglades is the only place in the world where both crocodiles and alligators coexist. The saltwater crocodile inhabits estuaries and mangroves. Snapping turtles also inhabit wetlands.
Birds, particularly waterfowl and waders use wetlands extensively.
Mammals of wetlands include numerous small and medium-sized species such as voles, bats, muskrats and platypus in addition to large herbivorous and apex predator species such as the beavers, coypu, swamp rabbit, Florida panther, jaguar, and moose. Wetlands attract many mammals due to abundant seeds, berries, and other vegetation as food for herbivores, as well as abundant populations of invertebrates, small reptiles and amphibians as prey for predators.
Invertebrates of wetlands include aquatic insects such as dragonflies, aquatic bugs and beetles, midges, mosquitos, crustaceans such as crabs, crayfish, shrimps, microcrustaceans, mollusks like clams, mussels, snails and worms. Invertebrates comprise more than half of the known animal species in wetlands, and are considered the primary food web link between plants and higher animals (such as fish and birds).
Ecosystem services
Depending on a wetland's geographic and topographic location, the functions it performs can support multiple ecosystem services, values, or benefits. United Nations Millennium Ecosystem Assessment and Ramsar Convention described wetlands as a whole to be of biosphere significance and societal importance in the following areas:
Water storage (flood control)
Groundwater replenishment
Shoreline stabilization and storm protection
Water purification
Wastewater treatment (in constructed wetlands)
Reservoirs of biodiversity
Pollination
Wetland products
Cultural values
Recreation and tourism
Climate change mitigation and adaptation
According to the Ramsar Convention:
The economic worth of the ecosystem services provided to society by intact, naturally functioning wetlands is frequently much greater than the perceived benefits of converting them to 'more valuable' intensive land use – particularly as the profits from unsustainable use often go to relatively few individuals or corporations, rather than being shared by society as a whole.
To replace these wetland ecosystem services, enormous amounts of money would need to be spent on water purification plants, dams, levees, and other hard infrastructure, and many of the services are impossible to replace.
Storage reservoirs and flood protection
Floodplains and closed-depression wetlands can provide the functions of storage reservoirs and flood protection. The wetland system of floodplains is formed from major rivers downstream from their headwaters. "The floodplains of major rivers act as natural storage reservoirs, enabling excess water to spread out over a wide area, which reduces its depth and speed. Wetlands close to the headwaters of streams and rivers can slow down rainwater runoff and spring snowmelt so that it does not run straight off the land into water courses. This can help prevent sudden, damaging floods downstream."
Notable river systems that produce wide floodplains include the Nile River, the Niger river inland delta, the Zambezi River flood plain, the Okavango River inland delta, the Kafue River flood plain, the Lake Bangweulu flood plain (Africa), Mississippi River (US), Amazon River (South America), Yangtze River (China), Danube River (Central Europe) and Murray-Darling River (Australia).
Groundwater replenishment
Groundwater replenishment can be achieved for example by marsh, swamp, and subterranean karst and cave hydrological systems. The surface water visibly seen in wetlands only represents a portion of the overall water cycle, which also includes atmospheric water (precipitation) and groundwater. Many wetlands are directly linked to groundwater and they can be a crucial regulator of both the quantity and quality of water found below the ground. Wetlands that have permeable substrates like limestone or occur in areas with highly variable and fluctuating water tables have especially important roles in groundwater replenishment or water recharge.
Substrates that are porous allow water to filter down through the soil and underlying rock into aquifers which are the source of much of the world's drinking water. Wetlands can also act as recharge areas when the surrounding water table is low and as a discharge zone when it is high.
Shoreline stabilization and storm protection
Mangroves, coral reefs, salt marsh can help with shoreline stabilization and storm protection. Tidal and inter-tidal wetland systems protect and stabilize coastal zones. Coral reefs provide a protective barrier to coastal shoreline. Mangroves stabilize the coastal zone from the interior and will migrate with the shoreline to remain adjacent to the boundary of the water. The main conservation benefit these systems have against storms and storm surges is the ability to reduce the speed and height of waves and floodwaters.
The United Kingdom has begun the concept of managed coastal realignment. This management technique provides shoreline protection through restoration of natural wetlands rather than through applied engineering. In East Asia, reclamation of coastal wetlands has resulted in widespread transformation of the coastal zone, and up to 65% of coastal wetlands have been destroyed by coastal development. One analysis using the impact of hurricanes versus storm protection provided naturally by wetlands projected the value of this service at US$33,000/hectare/year.
Water purification
Water purification can be provided by floodplains, closed-depression wetlands, mudflat, freshwater marsh, salt marsh, mangroves. Wetlands cycle both sediments and nutrients, sometimes serving as buffers between terrestrial and aquatic ecosystems. A natural function of wetland vegetation is the up-take, storage, and (for nitrate) the removal of nutrients found in runoff water from the surrounding landscapes.
Precipitation and surface runoff induces soil erosion, transporting sediment in suspension into and through waterways. All types of sediments whether composed of clay, silt, sand or gravel and rock can be carried into wetland systems through erosion. Wetland vegetation acts as a physical barrier to slow water flow and then trap sediment for both short or long periods of time. Suspended sediment can contain heavy metals that are also retained when wetlands trap the sediment.
The ability of wetland systems to store or remove nutrients and trap sediment is highly efficient and effective but each system has a threshold. An overabundance of nutrient input from fertilizer run-off, sewage effluent, or non-point pollution will cause eutrophication. Upstream erosion from deforestation can overwhelm wetlands making them shrink in size and cause dramatic biodiversity loss through excessive sedimentation load.
Wastewater treatment
Constructed wetlands are built for wastewater treatment. An example of how a natural wetland is used to provide some degree of sewage treatment is the East Kolkata Wetlands in Kolkata, India. The wetlands cover , and are used to treat Kolkata's sewage. The nutrients contained in the wastewater sustain fish farms and agriculture.
Reservoirs of biodiversity
Wetland systems' rich biodiversity has become a focal point catalysed by the Ramsar Convention and World Wildlife Fund. The impact of maintaining biodiversity is seen at the local level through job creation, sustainability, and community productivity. A good example is the Lower Mekong basin which runs through Cambodia, Laos, and Vietnam, supporting over 55 million people.
A key fish species which is overfished, the Piramutaba catfish, Brachyplatystoma vaillantii, migrates more than from its nursery grounds near the mouth of the Amazon River to its spawning grounds in Andean tributaries, above sea level, distributing plant seeds along the route.
Intertidal mudflats have a level of productivity similar to that of some wetlands even while possessing a low number of species. The abundant invertebrates found within the mud are a food source for migratory waterfowl.
Mudflats, saltmarshes, mangroves, and seagrass beds have high levels of both species richness and productivity, and are home to important nursery areas for many commercial fish stocks.
Populations of many species are confined geographically to only one or a few wetland systems, often due to the long period of time that the wetlands have been physically isolated from other aquatic sources. For example, the number of endemic species in the Selenga River Delta of Lake Baikal in Russia classifies it as a hotspot for biodiversity and one of the most biodiverse wetlands in the entire world.
Wetland products
Wetlands naturally produce an array of vegetation and other ecological products that can be harvested for personal and commercial use. Many fishes have all or part of their life-cycle occurring within a wetland system. Fresh and saltwater fish are the main source of protein for about one billion people and comprise 15% of an additional 3.5 billion people's protein intake. Another food staple found in wetland systems is rice, a popular grain that is consumed at the rate of one fifth of the total global calorie count. In Bangladesh, Cambodia and Vietnam, where rice paddies are predominant on the landscape, rice consumption reach 70%. Some native wetland plants in the Caribbean and Australia are harvested sustainably for medicinal compounds; these include the red mangrove (Rhizophora mangle) which possesses antibacterial, wound-healing, anti-ulcer effects, and antioxidant properties.
Other mangrove-derived products include fuelwood, salt (produced by evaporating seawater), animal fodder, traditional medicines (e.g. from mangrove bark), fibers for textiles and dyes and tannins.
Additional services and uses of wetlands
Some types of wetlands can serve as fire breaks that help slow the spread of minor wildfires. Larger wetland systems can influence local precipitation patterns. Some boreal wetland systems in catchment headwaters may help extend the period of flow and maintain water temperature in connected downstream waters. Pollination services are supported by many wetlands which may provide the only suitable habitat for pollinating insects, birds, and mammals in highly developed areas.
Disturbances and human impacts
Wetlands, the functions and services they provide as well as their flora and fauna, can be affected by several types of disturbances. The disturbances (sometimes termed stressors or alterations) can be human-associated or natural, direct or indirect, reversible or not, and isolated or cumulative.
Disturbances include exogenous factors such as flooding or drought. Humans are disturbing and damaging wetlands for example by oil and gas extraction, building infrastructure, overgrazing of livestock, overfishing, alteration of wetlands including dredging and draining, nutrient pollution and water pollution. Disturbance puts different levels of stress on an environment depending on the type and duration of disturbance.
Predominant disturbances of wetlands include:
Enrichment/eutrophication
Organic loading and reduced dissolved oxygen
Contaminant toxicity
Acidification
Salinization
Sedimentation
Altered solar input (turbidity/shade)
Vegetation removal
Thermal alteration
Drying/aridification
Inundation/flooding
Habitat fragmentation
Other human impacts
Disturbances can be further categorized as follows:
Minor disturbance: Stress that maintains ecosystem integrity.
Moderate disturbance: Ecosystem integrity is damaged but can recover in time without assistance.
Impairment or severe disturbance: Human intervention may be needed in order for ecosystem to recover.
Nutrient pollution comes from nitrogen inputs to aquatic systems and have drastically effected the dissolved nitrogen content of wetlands, introducing higher nutrient availability which leads to eutrophication.
Biodiversity loss occurs in wetland systems through land use changes, habitat destruction, pollution, exploitation of resources, and invasive species. For example, the introduction of water hyacinth, a native plant of South America into Lake Victoria in East Africa as well as duckweed into non-native areas of Queensland, Australia, have overtaken entire wetland systems overwhelming the habitats and reducing the diversity of native plants and animals.
Conversion to dry land
To increase economic productivity, wetlands are often converted into dry land with dykes and drains and used for agricultural purposes. The construction of dykes, and dams, has negative consequences for individual wetlands and entire watersheds. Their proximity to lakes and rivers means that they are often developed for human settlement. Once settlements are constructed and protected by dykes, the settlements then become vulnerable to land subsidence and ever increasing risk of flooding. The Mississippi River Delta around New Orleans, Louisiana is a well-known example; the Danube Delta in Europe is another.
Drainage of floodplains
Drainage of floodplains or development activities that narrow floodplain corridors (such as the construction of levees) reduces the ability of coupled river-floodplain systems to control flood damage. That is because modified and less expansive systems must still manage the same amount of precipitation, causing flood peaks to be higher or deeper and floodwaters to travel faster.
Water management engineering developments in the past century have degraded floodplain wetlands through the construction of artificial embankments such as dykes, bunds, levees, weirs, barrages and dams. All concentrate water into a main channel and waters that historically spread slowly over a large, shallow area are concentrated. Loss of wetland floodplains results in more severe and damaging flooding. Catastrophic human impact in the Mississippi River floodplains was seen in death of several hundred individuals during a levee breach in New Orleans caused by Hurricane Katrina. Human-made embankments along the Yangtze River floodplains have caused the main channel of the river to become prone to more frequent and damaging flooding. Some of these events include the loss of riparian vegetation, a 30% loss of the vegetation cover throughout the river's basin, a doubling of the percentage of the land affected by soil erosion, and a reduction in reservoir capacity through siltation build-up in floodplain lakes.
Overfishing
Overfishing is a major problem for sustainable use of wetlands. Concerns are developing over certain aspects of farm fishing, which uses natural wetlands and waterways to harvest fish for human consumption. Aquaculture is continuing to develop rapidly throughout the Asia-Pacific region especially in China where 90% of the total number of aquaculture farms occur, contributing 80% of global value. Some aquaculture has eliminated massive areas of wetland through practices such as the shrimp farming industry's destruction of mangroves. Even though the damaging impact of large-scale shrimp farming on the coastal ecosystem in many Asian countries has been widely recognized for quite some time now, it has proved difficult to mitigate since other employment avenues for people are lacking. Also burgeoning demand for shrimp globally has provided a large and ready market.
Conservation
Wetlands have historically subjected to large draining efforts for development (real estate or agriculture), and flooding to create recreational lakes or generate hydropower. Some of the world's most important agricultural areas were wetlands that have been converted to farmland. Since the 1970s, more focus has been put on preserving wetlands for their natural functions. Since 1900, 65–70% of the world's wetlands have been lost. In order to maintain wetlands and sustain their functions, alterations and disturbances that are outside the normal range of variation should be minimized.
Balancing wetland conservation with the needs of people
Wetlands are vital ecosystems that enhance the livelihoods for the millions of people who live in and around them. Studies have shown that it is possible to conserve wetlands while improving the livelihoods of people living among them. Case studies conducted in Malawi and Zambia looked at how dambos – wet, grassy valleys or depressions where water seeps to the surface – can be farmed sustainably. Project outcomes included a high yield of crops, development of sustainable farming techniques, and water management strategies that generate enough water for irrigation.
Ramsar Convention
The Ramsar Convention (full title: Convention on Wetlands of International Importance, especially as Waterfowl Habitat), is an international treaty designed to address global concerns regarding wetland loss and degradation. The primary purposes of the treaty are to list wetlands of international importance and to promote their wise use, with the ultimate goal of preserving the world's wetlands. Methods include restricting access to some wetland areas, as well as educating the public to combat the misconception that wetlands are wastelands. The Convention works closely with five International Organisation Partners (IOPs). These are: Birdlife International, the IUCN, the International Water Management Institute, Wetlands International and the World Wide Fund for Nature. The partners provide technical expertise, help conduct or facilitate field studies and provide financial support.
Restoration
Restoration and restoration ecologists intend to return wetlands to their natural trajectory by aiding directly with the natural processes of the ecosystem. These direct methods vary with respect to the degree of physical manipulation of the natural environment and each are associated with different levels of restoration. Restoration is needed after disturbance or perturbation of a wetland. There is no one way to restore a wetland and the level of restoration required will be based on the level of disturbance although, each method of restoration does require preparation and administration.
Levels of restoration
Factors influencing selected approach may include budget, time scale limitations, project goals, level of disturbance, landscape and ecological constraints, political and administrative agendas and socioeconomic priorities.
Prescribed natural or assisted regeneration
For this strategy, there is no biophysical manipulation and the ecosystem is left to recover based on the process of succession alone. The focus is to eliminate and prevent further disturbance from occurring and for this type of restoration requires prior research to understand the probability that the wetland will recover naturally. This is likely to be the first method of approach since it is the least intrusive and least expensive although some biophysical non-intrusive manipulation may be required to enhance the rate of succession to an acceptable level. Example methods include prescribed burns to small areas, promotion of site specific soil microbiota and plant growth using nucleation planting whereby plants radiate from an initial planting site, and promotion of niche diversity or increasing the range of niches to promote use by a variety of different species. These methods can make it easier for the natural species to flourish by removing environmental impediments and can speed up the process of succession.
Partial reconstruction
For this strategy, a mixture of natural regeneration and manipulated environmental control is used. This may require some engineering, and more intensive biophysical manipulations including ripping of subsoil, agrichemical applications of herbicides or insecticides, laying of mulch, mechanical seed dispersal, and tree planting on a large scale. In these circumstances the wetland is impaired and without human assistance it would not recover within an acceptable period of time as determined by ecologists. Methods of restoration used will have to be determined on a site by site basis as each location will require a different approach based on levels of disturbance and the local ecosystem dynamics.
Complete reconstruction
This most expensive and intrusive method of reconstruction requires engineering and ground up reconstruction. Because there is a redesign of the entire ecosystem it is important that the natural trajectory of the ecosystem be considered and that the plant species promoted will eventually return the ecosystem towards its natural trajectory.
In many cases constructed wetlands are often designed to treat stormwater/wastewater runoff. They can be used in developments as part of water-sensitive urban design systems and have benefits such as flood mitigation, removing pollutants, carbon sequestration, providing habitat for wildlife and biodiversity in often highly urbanised and fragmented landscapes.
Traditional knowledge
The ideas from traditional ecological knowledge can be applied as a holistic approach to the restoration of wetlands. These ideas focus more on responding to the observations detected from the environment considering that each part of a wetland ecosystem is interconnected. Applying these practices on specific locations of wetlands increase productivity, biodiversity, and improve its resilience. These practices include monitoring wetland resources, planting propagules, and addition of key species in order to create a self-sustaining wetland ecosystem.
Climate change aspects
Greenhouse gas emissions
In Southeast Asia, peat swamp forests and soils are being drained, burnt, mined, and overgrazed, contributing to climate change. As a result of peat drainage, the organic carbon that had built up over thousands of years and is normally under water is suddenly exposed to the air. The peat decomposes and is converted into carbon dioxide (CO2), which is then released into the atmosphere. Peat fires cause the same process to occur rapidly and in addition create enormous clouds of smoke that cross international borders, which now happens almost yearly in Southeast Asia. While peatlands constitute only 3% of the world's land area, their degradation produces 7% of all CO2 emissions.
Climate change mitigation
Studies have favorably identified the potential for coastal wetlands (also called blue carbon ecosystems) to provide some degree of climate change mitigation in two ways: by conservation, reducing the greenhouse gas emissions arising from the loss and degradation of such habitats, and by restoration, to increase carbon dioxide drawdown and its long-term storage. However, CO2 removal using coastal blue carbon restoration has questionable cost-effectiveness when considered only as a climate mitigation action, either for carbon-offsetting or for inclusion in Nationally Determined Contributions.
When wetlands are restored they have mitigation effects through their ability to sink carbon, converting a greenhouse gas (carbon dioxide) to solid plant material through the process of photosynthesis, and also through their ability to store and regulate water.
Wetlands store approximately 44.6 million tonnes of carbon per year globally (estimate from 2003). In salt marshes and mangrove swamps in particular, the average carbon sequestration rate is while peatlands sequester approximately .
Coastal wetlands, such as tropical mangroves and some temperate salt marshes, are known to be sinks for carbon that otherwise contribute to climate change in its gaseous forms (carbon dioxide and methane). The ability of many tidal wetlands to store carbon and minimize methane flux from tidal sediments has led to sponsorship of blue carbon initiatives that are intended to enhance those processes.
Climate change adaptation
The restoration of coastal blue carbon ecosystems is highly advantageous for climate change adaptation, coastal protection, food provision and biodiversity conservation.
Since the middle of the 20th century, human-caused climate change has resulted in observable changes in the global water cycle. A warming climate makes extremely wet and very dry occurrences more severe, causing more severe floods and droughts. For this reason, some of the ecosystem services that wetlands provide (e.g. water storage and flood control, groundwater replenishment, shoreline stabilization and storm protection) are important for climate change adaptation measures. In most parts of the world and under all emission scenarios, water cycle variability and accompanying extremes are anticipated to rise more quickly than the changes of average values.
Valuation
The value of a wetland to local communities typically involves first mapping a region's wetlands, then assessing the functions and ecosystem services the wetlands provide individually and cumulatively, and finally evaluating that information to prioritize or rank individual wetlands or wetland types for conservation, management, restoration, or development. Over the longer term, it requires keeping inventories of known wetlands and monitoring a representative sample of the wetlands to determine changes due to both natural and human factors.
Assessment
Rapid assessment methods are used to score, rank, rate, or categorize various functions, ecosystem services, species, communities, levels of disturbance, and/or ecological health of a wetland or group of wetlands. This is often done to prioritize particular wetlands for conservation (avoidance) or to determine the degree to which loss or alteration of wetland functions should be compensated, such as by restoring degraded wetlands elsewhere or providing additional protections to existing wetlands. Rapid assessment methods are also applied before and after a wetland has been restored or altered, to help monitor or predict the effects of those actions on various wetland functions and the services they provide. Assessments are typically considered to be "rapid" when they require only a single visit to the wetland lasting less than one day, which in some cases may include interpretation of aerial imagery and geographic information system (GIS) analyses of existing spatial data, but not detailed post-visit laboratory analyses of water or biological samples.
To achieve consistency among persons doing the assessment, rapid methods present indicator variables as questions or checklists on standardized data forms, and most methods standardize the scoring or rating procedure that is used to combine question responses into estimates of the levels of specified functions relative to the levels estimated in other wetlands ("calibration sites") assessed previously in a region. Rapid assessment methods, partly because they often use dozens of indicators pertaining to conditions surrounding a wetland as well as within the wetland itself, aim to provide estimates of wetland functions and services that are more accurate and repeatable than simply describing a wetland's class type. A need for wetland assessments to be rapid arises mostly when government agencies set deadlines for decisions affecting a wetland, or when the number of wetlands needing information on their functions or condition is large.
Inventory
Although developing a global inventory of wetlands has proven to be a large and difficult undertaking, many efforts at more local scales have been successful. Current efforts are based on available data, but both classification and spatial resolution have sometimes proven to be inadequate for regional or site-specific environmental management decision making. It is difficult to identify small, long, and narrow wetlands within the landscape. Many of today's remote sensing satellites do not have sufficient spatial and spectral resolution to monitor wetland conditions, although multispectral IKONOS and QuickBird data may offer improved spatial resolutions once it is 4 m or higher. Majority of the pixels are just mixtures of several plant species or vegetation types and are difficult to isolate which translates into an inability to classify the vegetation that defines the wetland. The growing availability of 3D vegetation and topography data from LiDAR has partially addressed the limitation of traditional multispectral imagery, as demonstrated in some case studies across the world.
Monitoring and mapping
A wetland needs to be monitored over time to assess whether it is functioning at an ecologically sustainable level or whether it is becoming degraded. Degraded wetlands will suffer a loss in water quality, loss of sensitive species, and aberrant functioning of soil geochemical processes.
Practically, many natural wetlands are difficult to monitor from the ground as they quite often are difficult to access and may require exposure to dangerous plants and animals as well as diseases borne by insects or other invertebrates. Remote sensing such as aerial imagery and satellite imaging provides effective tools to map and monitor wetlands across large geographic regions and over time. Many remote sensing methods can be used to map wetlands. The integration of multi-sourced data such as LiDAR and aerial photos proves more effective at mapping wetlands than the use of aerial photos alone, especially with the aid of modern machine learning methods (e.g., deep learning). Overall, using digital data provides a standardized data-collection procedure and an opportunity for data integration within a geographic information system.
Legislation
International efforts
National efforts
United States
Each country and region tends to have a codified definition of wetlands for legal purposes. In the United States, wetlands are defined as "those areas that are inundated or saturated by surface or groundwater at a frequency and duration sufficient to support, and that under normal circumstances do support, a prevalence of vegetation typically adapted for life in saturated soil conditions. Wetlands generally include swamps, marshes, bogs and similar areas". This definition has been used in the enforcement of the Clean Water Act. Some US states, such as Massachusetts and New York, have separate definitions that may differ from the federal government's.
In the United States Code, the term wetland is defined "as land that (A) has a predominance of hydric soils, (B) is inundated or saturated by surface or groundwater at a frequency and duration sufficient to support a prevalence of hydrophytic vegetation typically adapted for life in saturated soil conditions and (C) under normal circumstances supports a prevalence of such vegetation." Related to these legal definitions, "normal circumstances" are expected to occur during the wet portion of the growing season under normal climatic conditions (not unusually dry or unusually wet) and in the absence of significant disturbance. It is not uncommon for a wetland to be dry for long portions of the growing season. Still, under normal environmental conditions, the soils will be inundated to the surface, creating anaerobic conditions persisting through the wet portion of the growing season.
Canada
Wetlands and wetland policies in Canada
Other Individual Provincial and Territorial Based Policies
Examples
The world's largest wetlands include the swamp forests of the Amazon River basin, the peatlands of the West Siberian Plain, the Pantanal in South America, and the Sundarbans in the Ganges-Brahmaputra delta.
| Physical sciences | Wetlands | null |
102140 | https://en.wikipedia.org/wiki/Perturbation%20theory | Perturbation theory | In mathematics and applied mathematics, perturbation theory comprises methods for finding an approximate solution to a problem, by starting from the exact solution of a related, simpler problem. A critical feature of the technique is a middle step that breaks the problem into "solvable" and "perturbative" parts. In regular perturbation theory, the solution is expressed as a power series in a small parameter The first term is the known solution to the solvable problem. Successive terms in the series at higher powers of usually become smaller. An approximate 'perturbation solution' is obtained by truncating the series, often keeping only the first two terms, the solution to the known problem and the 'first order' perturbation correction.
Perturbation theory is used in a wide range of fields and reaches its most sophisticated and advanced forms in quantum field theory. Perturbation theory (quantum mechanics) describes the use of this method in quantum mechanics. The field in general remains actively and heavily researched across multiple disciplines.
Description
Perturbation theory develops an expression for the desired solution in terms of a formal power series known as a perturbation series in some "small" parameter, that quantifies the deviation from the exactly solvable problem. The leading term in this power series is the solution of the exactly solvable problem, while further terms describe the deviation in the solution, due to the deviation from the initial problem. Formally, we have for the approximation to the full solution a series in the small parameter (here called ), like the following:
In this example, would be the known solution to the exactly solvable initial problem, and the terms represent the first-order, second-order, third-order, and higher-order terms, which may be found iteratively by a mechanistic but increasingly difficult procedure. For small these higher-order terms in the series generally (but not always) become successively smaller. An approximate "perturbative solution" is obtained by truncating the series, often by keeping only the first two terms, expressing the final solution as a sum of the initial (exact) solution and the "first-order" perturbative correction
Some authors use big O notation to indicate the order of the error in the approximate solution:
If the power series in converges with a nonzero radius of convergence, the perturbation problem is called a regular perturbation problem. In regular perturbation problems, the asymptotic solution smoothly approaches the exact solution. However, the perturbation series can also diverge, and the truncated series can still be a good approximation to the true solution if it is truncated at a point at which its elements are minimum. This is called an asymptotic series. If the perturbation series is divergent or not a power series (for example, if the asymptotic expansion must include non-integer powers or negative powers ) then the perturbation problem is called a singular perturbation problem. Many special techniques in perturbation theory have been developed to analyze singular perturbation problems.
Prototypical example
The earliest use of what would now be called perturbation theory was to deal with the otherwise unsolvable mathematical problems of celestial mechanics: for example the orbit of the Moon, which moves noticeably differently from a simple Keplerian ellipse because of the competing gravitation of the Earth and the Sun.
Perturbation methods start with a simplified form of the original problem, which is simple enough to be solved exactly. In celestial mechanics, this is usually a Keplerian ellipse. Under Newtonian gravity, an ellipse is exactly correct when there are only two gravitating bodies (say, the Earth and the Moon) but not quite correct when there are three or more objects (say, the Earth, Moon, Sun, and the rest of the Solar System) and not quite correct when the gravitational interaction is stated using formulations from general relativity.
Perturbative expansion
Keeping the above example in mind, one follows a general recipe to obtain the perturbation series. The perturbative expansion is created by adding successive corrections to the simplified problem. The corrections are obtained by forcing consistency between the unperturbed solution, and the equations describing the system in full. Write for this collection of equations; that is, let the symbol stand in for the problem to be solved. Quite often, these are differential equations, thus, the letter "D".
The process is generally mechanical, if laborious. One begins by writing the equations so that they split into two parts: some collection of equations which can be solved exactly, and some additional remaining part for some small The solution (to ) is known, and one seeks the general solution to
Next the approximation is inserted into . This results in an equation for which, in the general case, can be written in closed form as a sum over integrals over Thus, one has obtained the first-order correction and thus is a good approximation to It is a good approximation, precisely because the parts that were ignored were of size The process can then be repeated, to obtain corrections and so on.
In practice, this process rapidly explodes into a profusion of terms, which become extremely hard to manage by hand. Isaac Newton is reported to have said, regarding the problem of the Moon's orbit, that "It causeth my head to ache." This unmanageability has forced perturbation theory to develop into a high art of managing and writing out these higher order terms. One of the fundamental breakthroughs in quantum mechanics for controlling the expansion are the Feynman diagrams, which allow quantum mechanical perturbation series to be represented by a sketch.
Examples
Perturbation theory has been used in a large number of different settings in physics and applied mathematics. Examples of the "collection of equations" include algebraic equations,
differential equations (e.g., the equations of motion
and commonly wave equations), thermodynamic free energy in statistical mechanics, radiative transfer,
and Hamiltonian operators in quantum mechanics.
Examples of the kinds of solutions that are found perturbatively include the solution of the equation of motion (e.g., the trajectory of a particle), the statistical average of some physical quantity (e.g., average magnetization), and the ground state energy of a quantum mechanical problem.
Examples of exactly solvable problems that can be used as starting points include linear equations, including linear equations of motion (harmonic oscillator, linear wave equation), statistical or quantum-mechanical systems of non-interacting particles (or in general, Hamiltonians or free energies containing only terms quadratic in all degrees of freedom).
Examples of systems that can be solved with perturbations include systems with nonlinear contributions to the equations of motion, interactions between particles, terms of higher powers in the Hamiltonian/free energy.
For physical problems involving interactions between particles, the terms of the perturbation series may be displayed (and manipulated) using Feynman diagrams.
History
Perturbation theory was first devised to solve otherwise intractable problems in the calculation of the motions of planets in the solar system. For instance, Newton's law of universal gravitation explained the gravitation between two astronomical bodies, but when a third body is added, the problem was, "How does each body pull on each?" Kepler's orbital equations only solve Newton's gravitational equations when the latter are limited to just two bodies interacting. The gradually increasing accuracy of astronomical observations led to incremental demands in the accuracy of solutions to Newton's gravitational equations, which led many eminent 18th and 19th century mathematicians, notably Joseph-Louis Lagrange and Pierre-Simon Laplace, to extend and generalize the methods of perturbation theory.
These well-developed perturbation methods were adopted and adapted to solve new problems arising during the development of quantum mechanics in 20th century atomic and subatomic physics. Paul Dirac developed quantum perturbation theory in 1927 to evaluate when a particle would be emitted in radioactive elements. This was later named Fermi's golden rule. Perturbation theory in quantum mechanics is fairly accessible, mainly because quantum mechanics is limited to linear wave equations, but also since the quantum mechanical notation allows expressions to be written in fairly compact form, thus making them easier to comprehend. This resulted in an explosion of applications, ranging from the Zeeman effect to the hyperfine splitting in the hydrogen atom.
Despite the simpler notation, perturbation theory applied to quantum field theory still easily gets out of hand. Richard Feynman developed the celebrated Feynman diagrams by observing that many terms repeat in a regular fashion. These terms can be replaced by dots, lines, squiggles and similar marks, each standing for a term, a denominator, an integral, and so on; thus complex integrals can be written as simple diagrams, with absolutely no ambiguity as to what they mean. The one-to-one correspondence between the diagrams, and specific integrals is what gives them their power. Although originally developed for quantum field theory, it turns out the diagrammatic technique is broadly applicable to many other perturbative series (although not always worthwhile).
In the second half of the 20th century, as chaos theory developed, it became clear that unperturbed systems were in general completely integrable systems, while the perturbed systems were not. This promptly lead to the study of "nearly integrable systems", of which the KAM torus is the canonical example. At the same time, it was also discovered that many (rather special) non-linear systems, which were previously approachable only through perturbation theory, are in fact completely integrable. This discovery was quite dramatic, as it allowed exact solutions to be given. This, in turn, helped clarify the meaning of the perturbative series, as one could now compare the results of the series to the exact solutions.
The improved understanding of dynamical systems coming from chaos theory helped shed light on what was termed the small denominator problem or small divisor problem. In the 19th century Poincaré observed (as perhaps had earlier mathematicians) that sometimes 2nd and higher order terms in the perturbative series have "small denominators": That is, they have the general form where and are some complicated expressions pertinent to the problem to be solved, and and are real numbers; very often they are the energy of normal modes. The small divisor problem arises when the difference is small, causing the perturbative correction to "blow up", becoming as large or maybe larger than the zeroth order term. This situation signals a breakdown of perturbation theory: It stops working at this point, and cannot be expanded or summed any further. In formal terms, the perturbative series is an asymptotic series: A useful approximation for a few terms, but at some point becomes less accurate if even more terms are added. The breakthrough from chaos theory was an explanation of why this happened: The small divisors occur whenever perturbation theory is applied to a chaotic system. The one signals the presence of the other.
Beginnings in the study of planetary motion
Since the planets are very remote from each other, and since their mass is small as compared to the mass of the Sun, the gravitational forces between the planets can be neglected, and the planetary motion is considered, to a first approximation, as taking place along Kepler's orbits, which are defined by the equations of the two-body problem, the two bodies being the planet and the Sun.
Since astronomic data came to be known with much greater accuracy, it became necessary to consider how the motion of a planet around the Sun is affected by other planets. This was the origin of the three-body problem; thus, in studying the system Moon-Earth-Sun, the mass ratio between the Moon and the Earth was chosen as the "small parameter". Lagrange and Laplace were the first to advance the view that the so-called "constants" which describe the motion of a planet around the Sun gradually change: They are "perturbed", as it were, by the motion of other planets and vary as a function of time; hence the name "perturbation theory".
Perturbation theory was investigated by the classical scholars – Laplace, Siméon Denis Poisson, Carl Friedrich Gauss – as a result of which the computations could be performed with a very high accuracy. The discovery of the planet Neptune in 1848 by Urbain Le Verrier, based on the deviations in motion of the planet Uranus. He sent the coordinates to J.G. Galle who successfully observed Neptune through his telescope – a triumph of perturbation theory.
Perturbation orders
The standard exposition of perturbation theory is given in terms of the order to which the perturbation is carried out: first-order perturbation theory or second-order perturbation theory, and whether the perturbed states are degenerate, which requires singular perturbation. In the singular case extra care must be taken, and the theory is slightly more elaborate.
In chemistry
Many of the ab initio quantum chemistry methods use perturbation theory directly or are closely related methods. Implicit perturbation theory works with the complete Hamiltonian from the very beginning and never specifies a perturbation operator as such. Møller–Plesset perturbation theory uses the difference between the Hartree–Fock Hamiltonian and the exact non-relativistic Hamiltonian as the perturbation. The zero-order energy is the sum of orbital energies. The first-order energy is the Hartree–Fock energy and electron correlation is included at second-order or higher. Calculations to second, third or fourth order are very common and the code is included in most ab initio quantum chemistry programs. A related but more accurate method is the coupled cluster method.
Shell-crossing
A shell-crossing (sc) occurs in perturbation theory when matter trajectories intersect, forming a singularity. This limits the predictive power of physical simulations at small scales.
| Physical sciences | Quantum mechanics | Physics |
102182 | https://en.wikipedia.org/wiki/Celestial%20mechanics | Celestial mechanics | Celestial mechanics is the branch of astronomy that deals with the motions of objects in outer space. Historically, celestial mechanics applies principles of physics (classical mechanics) to astronomical objects, such as stars and planets, to produce ephemeris data.
History
Modern analytic celestial mechanics started with Isaac Newton's Principia (1687). The name celestial mechanics is more recent than that. Newton wrote that the field should be called "rational mechanics". The term "dynamics" came in a little later with Gottfried Leibniz, and over a century after Newton, Pierre-Simon Laplace introduced the term celestial mechanics. Prior to Kepler, there was little connection between exact, quantitative prediction of planetary positions, using geometrical or numerical techniques, and contemporary discussions of the physical causes of the planets' motion.
Laws of planetary motion
Johannes Kepler as the first to closely integrate the predictive geometrical astronomy, which had been dominant from Ptolemy in the 2nd century to Copernicus, with physical concepts to produce a New Astronomy, Based upon Causes, or Celestial Physics in 1609. His work led to the laws of planetary orbits, which he developed using his physical principles and the planetary observations made by Tycho Brahe. Kepler's elliptical model greatly improved the accuracy of predictions of planetary motion, years before Newton developed his law of gravitation in 1686.
Newtonian mechanics and universal gravitation
Isaac Newton is credited with introducing the idea that the motion of objects in the heavens, such as planets, the Sun, and the Moon, and the motion of objects on the ground, like cannon balls and falling apples, could be described by the same set of physical laws. In this sense he unified celestial and terrestrial dynamics. Using his law of gravity, Newton confirmed Kepler's laws for elliptical orbits by deriving them from the gravitational two-body problem, which Newton included in his epochal Philosophiæ Naturalis Principia Mathematica in 1687.
Three-body problem
After Newton, Joseph-Louis Lagrange attempted to solve the three-body problem in 1772, analyzed the stability of planetary orbits, and discovered the existence of the Lagrange points. Lagrange also reformulated the principles of classical mechanics, emphasizing energy more than force, and developing a method to use a single polar coordinate equation to describe any orbit, even those that are parabolic and hyperbolic. This is useful for calculating the behaviour of planets and comets and such (parabolic and hyperbolic orbits are conic section extensions of Kepler's elliptical orbits). More recently, it has also become useful to calculate spacecraft trajectories.
Henri Poincaré published two now classical monographs, "New Methods of Celestial Mechanics" (1892–1899) and "Lectures on Celestial Mechanics" (1905–1910). In them, he successfully applied the results of their research to the problem of the motion of three bodies and studied in detail the behavior of solutions (frequency, stability, asymptotic, and so on). Poincaré showed that the three-body problem is not integrable. In other words, the general solution of the three-body problem can not be expressed in terms of algebraic and transcendental functions through unambiguous coordinates and velocities of the bodies. His work in this area was the first major achievement in celestial mechanics since Isaac Newton.
These monographs include an idea of Poincaré, which later became the basis for mathematical "chaos theory" (see, in particular, the Poincaré recurrence theorem) and the general theory of dynamical systems. He introduced the important concept of bifurcation points and proved the existence of equilibrium figures such as the non-ellipsoids, including ring-shaped and pear-shaped figures, and their stability. For this discovery, Poincaré received the Gold Medal of the Royal Astronomical Society (1900).
Standardisation of astronomical tables
Simon Newcomb was a Canadian-American astronomer who revised Peter Andreas Hansen's table of lunar positions. In 1877, assisted by George William Hill, he recalculated all the major astronomical constants. After 1884 he conceived, with A.M.W. Downing, a plan to resolve much international confusion on the subject. By the time he attended a standardisation conference in Paris, France, in May 1886, the international consensus was that all ephemerides should be based on Newcomb's calculations. A further conference as late as 1950 confirmed Newcomb's constants as the international standard.
Anomalous precession of Mercury
Albert Einstein explained the anomalous precession of Mercury's perihelion in his 1916 paper The Foundation of the General Theory of Relativity. General relativity led astronomers to recognize that Newtonian mechanics did not provide the highest accuracy.
Examples of problems
Celestial motion, without additional forces such as drag forces or the thrust of a rocket, is governed by the reciprocal gravitational acceleration between masses. A generalization is the n-body problem, where a number n of masses are mutually interacting via the gravitational force. Although analytically not integrable in the general case, the integration can be well approximated numerically.
Examples:
4-body problem: spaceflight to Mars (for parts of the flight the influence of one or two bodies is very small, so that there we have a 2- or 3-body problem; see also the patched conic approximation)
3-body problem:
Quasi-satellite
Spaceflight to, and stay at a Lagrangian point
In the case (two-body problem) the configuration is much simpler than for . In this case, the system is fully integrable and exact solutions can be found.
Examples:
A binary star, e.g., Alpha Centauri (approx. the same mass)
A binary asteroid, e.g., 90 Antiope (approx. the same mass)
A further simplification is based on the "standard assumptions in astrodynamics", which include that one body, the orbiting body, is much smaller than the other, the central body. This is also often approximately valid.
Examples:
The Solar System orbiting the center of the Milky Way
A planet orbiting the Sun
A moon orbiting a planet
A spacecraft orbiting Earth, a moon, or a planet (in the latter cases the approximation only applies after arrival at that orbit)
Perturbation theory
Perturbation theory comprises mathematical methods that are used to find an approximate solution to a problem which cannot be solved exactly. (It is closely related to methods used in numerical analysis, which are ancient.) The earliest use of modern perturbation theory was to deal with the otherwise unsolvable mathematical problems of celestial mechanics: Newton's solution for the orbit of the Moon, which moves noticeably differently from a simple Keplerian ellipse because of the competing gravitation of the Earth and the Sun.
Perturbation methods start with a simplified form of the original problem, which is carefully chosen to be exactly solvable. In celestial mechanics, this is usually a Keplerian ellipse, which is correct when there are only two gravitating bodies (say, the Earth and the Moon), or a circular orbit, which is only correct in special cases of two-body motion, but is often close enough for practical use.
The solved, but simplified problem is then "perturbed" to make its time-rate-of-change equations for the object's position closer to the values from the real problem, such as including the gravitational attraction of a third, more distant body (the Sun). The slight changes that result from the terms in the equations – which themselves may have been simplified yet again – are used as corrections to the original solution. Because simplifications are made at every step, the corrections are never perfect, but even one cycle of corrections often provides a remarkably better approximate solution to the real problem.
There is no requirement to stop at only one cycle of corrections. A partially corrected solution can be re-used as the new starting point for yet another cycle of perturbations and corrections. In principle, for most problems the recycling and refining of prior solutions to obtain a new generation of better solutions could continue indefinitely, to any desired finite degree of accuracy.
The common difficulty with the method is that the corrections usually progressively make the new solutions very much more complicated, so each cycle is much more difficult to manage than the previous cycle of corrections. Newton is reported to have said, regarding the problem of the Moon's orbit "It causeth my head to ache."
This general procedure – starting with a simplified problem and gradually adding corrections that make the starting point of the corrected problem closer to the real situation – is a widely used mathematical tool in advanced sciences and engineering. It is the natural extension of the "guess, check, and fix" method used anciently with numbers.
Reference frame
Problems in celestial mechanics are often posed in simplifying reference frames, such as the synodic reference frame applied to the three-body problem, where the origin coincides with the barycenter of the two larger celestial bodies. Other reference frames for n-body simulations include those that place the origin to follow the center of mass of a body, such as the heliocentric and the geocentric reference frames. The choice of reference frame gives rise to many phenomena, including the retrograde motion of superior planets while on a geocentric reference frame.
Orbital mechanics
| Physical sciences | Celestial mechanics | null |
102193 | https://en.wikipedia.org/wiki/Nonmetal | Nonmetal | In the context of the periodic table a nonmetal is a chemical element that mostly lacks distinctive metallic properties. They range from colorless gases like hydrogen to shiny crystals like iodine. Physically, they are usually lighter (less dense) than elements that form metals and are often poor conductors of heat and electricity. Chemically, nonmetals have relatively high electronegativity or usually attract electrons in a chemical bond with another element, and their oxides tend to be acidic.
Seventeen elements are widely recognized as nonmetals. Additionally, some or all of six borderline elements (metalloids) are sometimes counted as nonmetals.
The two lightest nonmetals, hydrogen and helium, together make up about 98% of the mass of the observable universe. Five nonmetallic elements—hydrogen, carbon, nitrogen, oxygen, and silicon—make up the bulk of Earth's atmosphere, biosphere, crust and oceans.
Industrial uses of nonmetals include in electronics, energy storage, agriculture, and chemical production.
Most nonmetallic elements were identified in the 18th and 19th centuries. While a distinction between metals and other minerals had existed since antiquity, a basic classification of chemical elements as metallic or nonmetallic emerged only in the late 18th century. Since then about twenty properties have been suggested as criteria for distinguishing nonmetals from metals.
Definition and applicable elements
Unless otherwise noted, this article describes the stable form of an element at standard temperature and pressure (STP).
Nonmetallic chemical elements are often described as lacking properties common to metals, namely shininess, pliability, good thermal and electrical conductivity, and a general capacity to form basic oxides. There is no widely accepted precise definition; any list of nonmetals is open to debate and revision. The elements included depend on the properties regarded as most representative of nonmetallic or metallic character.
Fourteen elements are almost always recognized as nonmetals:
Three more are commonly classed as nonmetals, but some sources list them as "metalloids", a term which refers to elements regarded as intermediate between metals and nonmetals:
One or more of the six elements most commonly recognized as metalloids are sometimes instead counted as nonmetals:
About 15–20% of the 118 known elements are thus classified as nonmetals.
General properties
Physical
Nonmetals vary greatly in appearance, being colorless, colored or shiny.
For the colorless nonmetals (hydrogen, nitrogen, oxygen, and the noble gases), no absorption of light happens in the visible part of the spectrum, and all visible light is transmitted.
The colored nonmetals (sulfur, fluorine, chlorine, bromine) absorb some colors (wavelengths) and transmit the complementary or opposite colors. For example, chlorine's "familiar yellow-green colour ... is due to a broad region of absorption in the violet and blue regions of the spectrum". The shininess of boron, graphite (carbon), silicon, black phosphorus, germanium, arsenic, selenium, antimony, tellurium, and iodine is a result of varying degrees of metallic conduction where the electrons can reflect incoming visible light.
About half of nonmetallic elements are gases under standard temperature and pressure; most of the rest are solids. Bromine, the only liquid, is usually topped by a layer of its reddish-brown fumes. The gaseous and liquid nonmetals have very low densities, melting and boiling points, and are poor conductors of heat and electricity. The solid nonmetals have low densities and low mechanical strength (being either hard and brittle, or soft and crumbly), and a wide range of electrical conductivity.
This diversity in form stems from variability in internal structures and bonding arrangements. Covalent nonmetals existing as discrete atoms like xenon, or as small molecules, such as oxygen, sulfur, and bromine, have low melting and boiling points; many are gases at room temperature, as they are held together by weak London dispersion forces acting between their atoms or molecules, although the molecules themselves have strong covalent bonds. In contrast, nonmetals that form extended structures, such as long chains of selenium atoms, sheets of carbon atoms in graphite, or three-dimensional lattices of silicon atoms have higher melting and boiling points, and are all solids, as it takes more energy to overcome their stronger bonding. Nonmetals closer to the left or bottom of the periodic table (and so closer to the metals) often have metallic interactions between their molecules, chains, or layers; this occurs in boron, carbon, phosphorus, arsenic, selenium, antimony, tellurium and iodine.
Covalently bonded nonmetals often share only the electrons required to achieve a noble gas electron configuration. For example, nitrogen forms diatomic molecules featuring a triple bonds between each atom, both of which thereby attain the configuration of the noble gas neon. Antimony's larger atomic size prevents triple bonding, resulting in buckled layers in which each antimony atom is singly bonded with three other nearby atoms.
Good electrical conductivity occurs when there is metallic bonding, however the electrons in nonmetals are often not metallic. Good electrical and thermal conductivity associated with metallic electrons is seen in carbon (as graphite, along its planes), arsenic, and antimony. Good thermal conductivity occurs in boron, silicon, phosphorus, and germanium; such conductivity is transmitted though vibrations of the crystalline lattices of these elements. Moderate electrical conductivity is observed in the semiconductors boron, silicon, phosphorus, germanium, selenium, tellurium, and iodine.
Many of the nonmetallic elements are hard and brittle, where dislocations cannot readily move so they tend to undergo brittle fracture rather than deforming. Some do deform such as white phosphorus (soft as wax, pliable and can be cut with a knife, at room temperature), in plastic sulfur, and in selenium which can be drawn into wires from its molten state. Graphite is a standard solid lubricant where dislocations move very easily in the basal planes.
Allotropes
Over half of the nonmetallic elements exhibit a range of less stable allotropic forms, each with distinct physical properties. For example, carbon, the most stable form of which is graphite, can manifest as diamond, buckminsterfullerene, amorphous and paracrystalline variations. Allotropes also occur for nitrogen, oxygen, phosphorus, sulfur, selenium and iodine.
Chemical
Nonmetals have relatively high values of electronegativity, and their oxides are usually acidic. Exceptions may occur if a nonmetal is not very electronegative, or if its oxidation state is low, or both. These non-acidic oxides of nonmetals may be amphoteric (like water, H2O) or neutral (like nitrous oxide, N2O), but never basic.
Nonmetals tend to gain electrons during chemical reactions, in contrast to metals which tend to donate electrons. This behavior is related to the stability of electron configurations in the noble gases, which have complete outer shells as summarized by the duet and octet rules of thumb, more correctly explained in terms of valence bond theory.
They typically exhibit higher ionization energies, electron affinities, and standard electrode potentials than metals. Generally, the higher these values are (including electronegativity) the more nonmetallic the element tends to be. For example, the chemically very active nonmetals fluorine, chlorine, bromine, and iodine have an average electronegativity of 3.19—a figure higher than that of any metallic element.
The chemical distinctions between metals and nonmetals is connected to the attractive force between the positive nuclear charge of an individual atom and its negatively charged outer electrons. From left to right across each period of the periodic table, the nuclear charge (number of protons in the atomic nucleus) increases. There is a corresponding reduction in atomic radius as the increased nuclear charge draws the outer electrons closer to the nuclear core. In chemical bonding, nonmetals tend to gain electrons due to their higher nuclear charge, resulting in negatively charged ions.
The number of compounds formed by nonmetals is vast. The first 10 places in a "top 20" table of elements most frequently encountered in 895,501,834 compounds, as listed in the Chemical Abstracts Service register for November 2, 2021, were occupied by nonmetals. Hydrogen, carbon, oxygen, and nitrogen collectively appeared in most (80%) of compounds. Silicon, a metalloid, ranked 11th. The highest-rated metal, with an occurrence frequency of 0.14%, was iron, in 12th place. A few examples of nonmetal compounds are: boric acid (), used in ceramic glazes; selenocysteine (), the 21st amino acid of life; phosphorus sesquisulfide (P4S3), found in strike anywhere matches; and teflon ()n), used to create non-stick coatings for pans and other cookware.
Complications
Adding complexity to the chemistry of the nonmetals are anomalies occurring in the first row of each periodic table block; non-uniform periodic trends; higher oxidation states; multiple bond formation; and property overlaps with metals.
First row anomaly
Starting with hydrogen, the first row anomaly primarily arises from the electron configurations of the elements concerned. Hydrogen is notable for its diverse bonding behaviors. It most commonly forms covalent bonds, but it can also lose its single electron in an aqueous solution, leaving behind a bare proton with tremendous polarizing power. Consequently, this proton can attach itself to the lone electron pair of an oxygen atom in a water molecule, laying the foundation for acid-base chemistry. Moreover, a hydrogen atom in a molecule can form a second, albeit weaker, bond with an atom or group of atoms in another molecule. Such bonding, "helps give snowflakes their hexagonal symmetry, binds DNA into a double helix; shapes the three-dimensional forms of proteins; and even raises water's boiling point high enough to make a decent cup of tea."
Hydrogen and helium, as well as boron through neon, have unusually small atomic radii. This phenomenon arises because the 1s and 2p subshells lack inner analogues (meaning there is no zero shell and no 1p subshell), and they therefore experience less electron-electron exchange interactions, unlike the 3p, 4p, and 5p subshells of heavier elements. As a result, ionization energies and electronegativities among these elements are higher than the periodic trends would otherwise suggest. The compact atomic radii of carbon, nitrogen, and oxygen facilitate the formation of double or triple bonds.
While it would normally be expected, on electron configuration consistency grounds, that hydrogen and helium would be placed atop the s-block elements, the significant first row anomaly shown by these two elements justifies alternative placements. Hydrogen is occasionally positioned above fluorine, in group 17, rather than above lithium in group 1. Helium is almost always placed above neon, in group 18, rather than above beryllium in group 2.
Secondary periodicity
An alternation in certain periodic trends, sometimes referred to as secondary periodicity, becomes evident when descending groups 13 to 15, and to a lesser extent, groups 16 and 17. Immediately after the first row of d-block metals, from scandium to zinc, the 3d electrons in the p-block elements—specifically, gallium (a metal), germanium, arsenic, selenium, and bromine—prove less effective at shielding the increasing positive nuclear charge.
The Soviet chemist gives two more tangible examples:
"The toxicity of some arsenic compounds, and the absence of this property in analogous compounds of phosphorus [P] and antimony [Sb]; and the ability of selenic acid [] to bring metallic gold [Au] into solution, and the absence of this property in sulfuric [] and [] acids."
Higher oxidation states
Roman numerals such as III, V and VIII denote oxidation states
Some nonmetallic elements exhibit oxidation states that deviate from those predicted by the octet rule, which typically results in an oxidation state of –3 in group 15, –2 in group 16, –1 in group 17, and 0 in group 18. Examples include ammonia NH3, hydrogen sulfide H2S, hydrogen fluoride HF, and elemental xenon Xe. Meanwhile, the maximum possible oxidation state increases from +5 in group 15, to +8 in group 18. The +5 oxidation state is observable from period 2 onward, in compounds such as nitric acid HN(V)O3 and phosphorus pentafluoride PCl5. Higher oxidation states in later groups emerge from period 3 onwards, as seen in sulfur hexafluoride SF6, iodine heptafluoride IF7, and xenon(VIII) tetroxide XeO4. For heavier nonmetals, their larger atomic radii and lower electronegativity values enable the formation of compounds with higher oxidation numbers, supporting higher bulk coordination numbers.
Multiple bond formation
Period 2 nonmetals, particularly carbon, nitrogen, and oxygen, show a propensity to form multiple bonds. The compounds formed by these elements often exhibit unique stoichiometries and structures, as seen in the various nitrogen oxides, which are not commonly found in elements from later periods.
Property overlaps
While certain elements have traditionally been classified as nonmetals and others as metals, some overlapping of properties occurs. Writing early in the twentieth century, by which time the era of modern chemistry had been well-established, Humphrey observed that:
... these two groups, however, are not marked off perfectly sharply from each other; some nonmetals resemble metals in certain of their properties, and some metals approximate in some ways to the non-metals.
Examples of metal-like properties occurring in nonmetallic elements include:
Silicon has an electronegativity (1.9) comparable with metals such as cobalt (1.88), copper (1.9), nickel (1.91) and silver (1.93);
The electrical conductivity of graphite exceeds that of some metals;
Selenium can be drawn into a wire;
Radon is the most metallic of the noble gases and begins to show some cationic behavior, which is unusual for a nonmetal; and
In extreme conditions, just over half of nonmetallic elements can form homopolyatomic cations.
Examples of nonmetal-like properties occurring in metals are:
Tungsten displays some nonmetallic properties, sometimes being brittle, having a high electronegativity, and forming only anions in aqueous solution, and predominately acidic oxides.
Gold, the "king of metals" has the highest electrode potential among metals, suggesting a preference for gaining rather than losing electrons. Gold's ionization energy is one of the highest among metals, and its electron affinity and electronegativity are high, with the latter exceeding that of some nonmetals. It forms the Au– auride anion and exhibits a tendency to bond to itself, behaviors which are unexpected for metals. In aurides (MAu, where M = Li–Cs), gold's behavior is similar to that of a halogen. Gold has a large enough nuclear potential that the electrons have to be considered with relativistic effects included which changes some of the properties.
A relatively recent development involves certain compounds of heavier p-block elements, such as silicon, phosphorus, germanium, arsenic and antimony, exhibiting behaviors typically associated with transition metal complexes. This is linked to a small energy gap between their filled and empty molecular orbitals, which are the regions in a molecule where electrons reside and where they can be available for chemical reactions. In such compounds, this allows for unusual reactivity with small molecules like hydrogen (H2), ammonia (NH3), and ethylene (C2H4), a characteristic previously observed primarily in transition metal compounds. These reactions may open new avenues in catalytic applications.
Types
Nonmetal classification schemes vary widely, with some accommodating as few as two subtypes and others identifying up to seven. For example, the periodic table in the Encyclopaedia Britannica recognizes noble gases, halogens, and other nonmetals, and splits the elements commonly recognized as metalloids between "other metals" and "other nonmetals". On the other hand, seven of twelve color categories on the Royal Society of Chemistry periodic table include nonmetals.
Starting on the right side of the periodic table, three types of nonmetals can be recognized:
the relatively inert noble gases—helium, neon, argon, krypton, xenon, radon;
the notably reactive halogen nonmetals—fluorine, chlorine, bromine, iodine; and
the mixed reactivity "unclassified nonmetals", a set with no widely used collective name—hydrogen, carbon, nitrogen, oxygen, phosphorus, sulfur, selenium. The descriptive phrase unclassified nonmetals is used here for convenience.
The elements in a fourth set are sometimes recognized as nonmetals:
the generally unreactive metalloids, sometimes considered a third category distinct from metals and nonmetals—boron, silicon, germanium, arsenic, antimony, tellurium.
The boundaries between these types are not sharp. Carbon, phosphorus, selenium, and iodine border the metalloids and show some metallic character, as does hydrogen.
The greatest discrepancy between authors occurs in metalloid "frontier territory". Some consider metalloids distinct from both metals and nonmetals, while others classify them as nonmetals. Some categorize certain metalloids as metals (e.g., arsenic and antimony due to their similarities to heavy metals). Metalloids resemble the elements universally considered "nonmetals" in having relatively low densities, high electronegativity, and similar chemical behavior.
Noble gases
Six nonmetals are classified as noble gases: helium, neon, argon, krypton, xenon, and the radioactive radon. In conventional periodic tables they occupy the rightmost column. They are called noble gases due to their exceptionally low chemical reactivity.
These elements exhibit similar properties, characterized by their colorlessness, odorlessness, and nonflammability. Due to their closed outer electron shells, noble gases possess weak interatomic forces of attraction, leading to exceptionally low melting and boiling points. As a consequence, they all exist as gases under standard conditions, even those with atomic masses surpassing many typically solid elements.
Chemically, the noble gases exhibit relatively high ionization energies, negligible or negative electron affinities, and high to very high electronegativities. The number of compounds formed by noble gases is in the hundreds and continues to expand, with most of these compounds involving the combination of oxygen or fluorine with either krypton, xenon, or radon.
Halogen nonmetals
While the halogen nonmetals are notably reactive and corrosive elements, they can also be found in everyday compounds like toothpaste (NaF); common table salt (NaCl); swimming pool disinfectant (NaBr); and food supplements (KI). The term "halogen" itself means "salt former".
Chemically, the halogen nonmetals exhibit high ionization energies, electron affinities, and electronegativity values, and are mostly relatively strong oxidizing agents. These characteristics contribute to their corrosive nature. All four elements tend to form primarily ionic compounds with metals, in contrast to the remaining nonmetals (except for oxygen) which tend to form primarily covalent compounds with metals. The highly reactive and strongly electronegative nature of the halogen nonmetals epitomizes nonmetallic character.
Unclassified nonmetals
Hydrogen behaves in some respects like a metallic element and in others like a nonmetal. Like a metallic element it can, for example, form a solvated cation in aqueous solution; it can substitute for alkali metals in compounds such as the chlorides (NaCl cf. HCl) and nitrates (KNO3 cf. HNO3), and in certain alkali metal complexes as a nonmetal. It attains this configuration by forming a covalent or ionic bond or, if it has initially given up its electron, by attaching itself to a lone pair of electrons.
Some or all of these nonmetals share several properties. Being generally less reactive than the halogens, most of them can occur naturally in the environment. They have significant roles in biology and geochemistry. Collectively, their physical and chemical characteristics can be described as "moderately non-metallic". Sometimes they have corrosive aspects. Carbon corrosion can occur in fuel cells. Untreated selenium in soils can lead to the formation of corrosive hydrogen selenide gas. Very different, when combined with metals, the unclassified nonmetals can form interstitial or refractory compounds due to their relatively small atomic radii and sufficiently low ionization energies. They also exhibit a tendency to bond to themselves, particularly in solid compounds. Additionally, diagonal periodic table relationships among these nonmetals mirror similar relationships among the metalloids.
Abundance, extraction, and uses
Abundance
The abundance of elements in the universe results from nuclear physics processes like nucleosynthesis and radioactive decay.
The volatile noble gas nonmetal elements are less abundant in the atmosphere than expected based their overall abundance due to cosmic nucleosynthesis. Mechanisms to explain this difference is an important aspect of planetary science. Even within that challenge, the nonmetal element is unexpectedly depleted. A possible explanation comes from theoretical models of the high pressures in the Earth's core suggest there may be around 1013 tons of xenon, in the form of stable XeFe3 and XeNi3 intermetallic compounds.
Five nonmetals—hydrogen, carbon, nitrogen, oxygen, and silicon—form the bulk of the directly observable structure of the Earth: about 73% of the crust, 93% of the biomass, 96% of the hydrosphere, and over 99% of the atmosphere, as shown in the accompanying table. Silicon and oxygen form highly stable tetrahedral structures, known as silicates. Here, "the powerful bond that unites the oxygen and silicon ions is the cement that holds the Earth's crust together."
In the biomass, the relative abundance of the first four nonmetals (and phosphorus, sulfur, and selenium marginally) is attributed to a combination of relatively small atomic size, and sufficient spare electrons. These two properties enable them to bind to one another and "some other elements, to produce a molecular soup sufficient to build a self-replicating system."
Extraction
Nine of the 23 nonmetallic elements are gases, or form compounds that are gases, and are extracted from natural gas or liquid air. These elements include hydrogen, helium, nitrogen, oxygen, neon, sulfur, argon, krypton, and xenon. For example, nitrogen and oxygen are extracted from air through fractional distillation of liquid air. This method capitalizes on their different boiling points to separate them efficiently. Sulfur was extracted using the Frasch process, which involved injecting superheated water into underground deposits to melt the sulfur, which is then pumped to the surface. This technique leveraged sulfur's low melting point relative to other geological materials. It is now obtained by reacting the hydrogen sulfide in natural gas, with oxygen. Water is formed, leaving the sulfur behind.
are extracted from the following sources:
Gases (3): hydrogen, from methane; helium, from natural gas; sulfur, from hydrogen sulfide in natural gas
Liquids (9): nitrogen, oxygen, neon, argon, krypton and xenon from liquid air; chlorine, bromine and iodine from brine
Solids (12): boron, from borates; carbon occurs naturally as graphite; silicon, from silica; phosphorus, from phosphates; iodine, from sodium iodate; radon, as a decay product from uranium ores; fluorine, from fluorite; germanium, arsenic, selenium, antimony and tellurium, from sulfides.
Uses
Uses of nonmetals and non-metallic elements are broadly categorized as domestic, industrial, attenuative (lubricative, retarding, insulating or cooling), and agricultural
Many have domestic and industrial applications in household accoutrements; medicine and pharmaceuticals; and lasers and lighting. They are components of mineral acids; and prevalent in plug-in hybrid vehicles; and smartphones.
A significant number have attenuative and agricultural applications. They are used in lubricants; and flame retardants and fire extinguishers. They can serve as inert air replacements; and are used in cryogenics and refrigerants. Their significance extends to agriculture, through their use in fertilizers.
Additionally, a smaller number of nonmetals or nonmetallic elements find specialized uses in explosives; and welding gases.
Taxonomical history
Background
Around 340 BCE, in Book III of his treatise Meteorology, the ancient Greek philosopher Aristotle categorized substances found within the Earth into metals and "fossiles". The latter category included various minerals such as realgar, ochre, ruddle, sulfur, cinnabar, and other substances that he referred to as "stones which cannot be melted".
Until the Middle Ages the classification of minerals remained largely unchanged, albeit with varying terminology. In the fourteenth century, the English alchemist Richardus Anglicus expanded upon the classification of minerals in his work Correctorium Alchemiae. In this text, he proposed the existence of two primary types of minerals. The first category, which he referred to as "major minerals", included well-known metals such as gold, silver, copper, tin, lead, and iron. The second category, labeled "minor minerals", encompassed substances like salts, atramenta (iron sulfate), alums, vitriol, arsenic, orpiment, sulfur, and similar substances that were not metallic bodies.
The term "nonmetallic" dates back to at least the 16th century. In his 1566 medical treatise, French physician Loys de L'Aunay distinguished substances from plant sources based on whether they originated from metallic or non-metallic soils.
Later, the French chemist Nicolas Lémery discussed metallic and nonmetallic minerals in his work Universal Treatise on Simple Drugs, Arranged Alphabetically published in 1699. In his writings, he contemplated whether the substance "cadmia" belonged to either the first category, akin to cobaltum (cobaltite), or the second category, exemplified by what was then known as calamine—a mixed ore containing zinc carbonate and silicate.
Organization of elements by types
Just as the ancients distinguished metals from other minerals, similar distinctions developed as the modern idea of chemical elements emerged in the late 1700s. French chemist Antoine Lavoisier published the first modern list of chemical elements in his revolutionary 1789 Traité élémentaire de chimie. The 33 elements known to Lavoisier were categorized into four distinct groups, including gases, metallic substances, nonmetallic substances that form acids when oxidized, and earths (heat-resistant oxides). Lavoisier's work gained widespread recognition and was republished in twenty-three editions across six languages within its first seventeen years, significantly advancing the understanding of chemistry in Europe and America.
In 1802 the term "metalloids" was introduced for elements with the physical properties of metals but the chemical properties of non-metals. However,
in 1811, the Swedish chemist Berzelius used the term "metalloids" to describe all nonmetallic elements, noting their ability to form negatively charged ions with oxygen in aqueous solutions.
Thus in 1864, the "Manual of Metalloids" divided all elements into either metals or metalloids, with the latter group including elements now called nonmetals. Reviews of the book indicated that the term "metalloids" was still endorsed by leading authorities, but there were reservations about its appropriateness. While Berzelius' terminology gained significant acceptance, it later faced criticism from some who found it counterintuitive, misapplied, or even invalid. The idea of designating elements like arsenic as metalloids had been considered. By as early as 1866, some authors began preferring the term "nonmetal" over "metalloid" to describe nonmetallic elements. In 1875, Kemshead observed that elements were categorized into two groups: non-metals (or metalloids) and metals. He noted that the term "non-metal", despite its compound nature, was more precise and had become universally accepted as the nomenclature of choice.
Development of types
In 1844, , a French doctor, pharmacist, and chemist, established a basic taxonomy of nonmetals to aid in their study. He wrote:
They will be divided into four groups or sections, as in the following:
Organogens—oxygen, nitrogen, hydrogen, carbon
Sulphuroids—sulfur, selenium, phosphorus
Chloroides—fluorine, chlorine, bromine, iodine
Boroids—boron, silicon.
Dupasquier's quartet parallels the modern nonmetal types. The organogens and sulphuroids are akin to the unclassified nonmetals. The chloroides were later called halogens. The boroids eventually evolved into the metalloids, with this classification beginning from as early as 1864. The then unknown noble gases were recognized as a distinct nonmetal group after being discovered in the late 1800s.
His taxonomy was noted for its natural basis. That said, it was a significant departure from other contemporary classifications, since it grouped together oxygen, nitrogen, hydrogen, and carbon.
In 1828 and 1859, the French chemist Dumas classified nonmetals as (1) hydrogen; (2) fluorine to iodine; (3) oxygen to sulfur; (4) nitrogen to arsenic; and (5) carbon, boron and silicon, thereby anticipating the vertical groupings of Mendeleev's 1871 periodic table. Dumas' five classes fall into modern groups 1, 17, 16, 15, and 14 to
13 respectively.
Suggested distinguishing criteria
Much of the early analyses were phenomenological, and a variety of physical, chemical, and atomic properties have been suggested for distinguishing metals from nonmetals (or other bodies); a comprehensive early set of characteristics was stated by Rev Thaddeus Mason Harrisn in the 1803 Minor Encyclopedia .
METAL, in natural history and chemistry, the name of a class of simple bodies; of which it is observed, that they posses; a lustre; that they are opaque; that they arc fusible, or may be melted; that their specific gravity is greater than that of any other bodies yet discovered; that they are better conductors of electricity, than any other body; that they are malleable, or capable of being extended and flattened by the hammer; and that they are ductile or tenacious, that is, capable of being drawn out into threads or wires.
Some criteria did not last long; for instance in 1809, the British chemist and inventor Humphry Davy isolated sodium and potassium, their low densities contrasted with their metallic appearance, so the density property was tenuous although these metals was firmly established by their chemical properties.
Johnson has a similar approach to Mason, distinguishing between metals and nonmetals on the basis of their physical states, electrical conductivity, mechanical properties, and the acid-base nature of their oxides:
gaseous elements are nonmetals (hydrogen, nitrogen, oxygen, fluorine, chlorine and the noble gases);
liquids (mercury, bromine) are either metallic or nonmetallic: mercury, as a good conductor, is a metal; bromine, with its poor conductivity, is a nonmetal;
solids are either ductile and malleable, hard and brittle, or soft and crumbly:
a. ductile and malleable elements are metals;
b. hard and brittle elements include boron, silicon and germanium, which are semiconductors and therefore not metals; and
c. soft and crumbly elements include carbon, phosphorus, sulfur, arsenic, antimony, tellurium and iodine, which have acidic oxides indicative of nonmetallic character.
Several authors have noted that nonmetals generally have low densities and high electronegativity. The accompanying table, using a threshold of 7 g/cm3 for density and 1.9 for electronegativity (revised Pauling), shows that all nonmetals have low density and high electronegativity. In contrast, all metals have either high density or low electronegativity (or both). Goldwhite and Spielman added that, "... lighter elements tend to be more electronegative than heavier ones." The average electronegativity for the elements in the table with densities less than 7 gm/cm3 (metals and nonmetals) is 1.97 compared to 1.66 for the metals having densities of more than 7 gm/cm3.
There is not full agreement about the use of phenomenological properties. Emsley pointed out the complexity of this task, asserting that no single property alone can unequivocally assign elements to either the metal or nonmetal category. Some authors divide elements into metals, metalloids, and nonmetals, but Oderberg disagrees, arguing that by the principles of categorization, anything not classified as a metal should be considered a nonmetal.
Kneen and colleagues proposed that the classification of nonmetals can be achieved by establishing a single criterion for metallicity. They acknowledged that various plausible classifications exist and emphasized that while these classifications may differ to some extent, they would generally agree on the categorization of nonmetals. The describe electrical conductivity as the key property, arguing that this is the most common approach.
One of the most commonly recognized properties used is the temperature coefficient of resistivity, the effect of heating on electrical resistance and conductivity. As temperature rises, the conductivity of metals decreases while that of nonmetals increases. However, plutonium, carbon, arsenic, and antimony appear to defy the norm. When plutonium (a metal) is heated within a temperature range of −175 to +125 °C its conductivity increases. Similarly, despite its common classification as a nonmetallic element, carbon (as graphite) is a semimetal which when heated experiences a decrease in electrical conductivity. Arsenic and antimony, which are occasionally classified as nonmetallic elements are also semimetals, and show behavior similar to carbon.
Comparison of selected properties
The two tables in this section list some of the properties of five types of elements (noble gases, halogen nonmetals, unclassified nonmetals, metalloids and, for comparison, metals) based on their most stable forms at standard temperature and pressure. The dashed lines around the columns for metalloids signify that the treatment of these elements as a distinct type can vary depending on the author, or classification scheme in use.
Physical properties by element type
Physical properties are listed in loose order of ease of their determination.
Chemical properties by element type
Chemical properties are listed from general characteristics to more specific details.
† Hydrogen can also form alloy-like hydrides
‡ The labels low, moderate, high, and very high are arbitrarily based on the value spans listed in the table
| Physical sciences | Periodic table | Chemistry |
102213 | https://en.wikipedia.org/wiki/Essential%20amino%20acid | Essential amino acid | An essential amino acid, or indispensable amino acid, is an amino acid that cannot be synthesized from scratch by the organism fast enough to supply its demand, and must therefore come from the diet. Of the 21 amino acids common to all life forms, the nine amino acids humans cannot synthesize are valine, isoleucine, leucine, methionine, phenylalanine, tryptophan, threonine, histidine, and lysine.
Six other amino acids are considered conditionally essential in the human diet, meaning their synthesis can be limited under special pathophysiological conditions, such as prematurity in the infant or individuals in severe catabolic distress. These six are arginine, cysteine, glycine, glutamine, proline, and tyrosine. Six amino acids are non-essential (dispensable) in humans, meaning they can be synthesized in sufficient quantities in the body. These six are alanine, aspartic acid, asparagine, glutamic acid, serine, and selenocysteine (considered the 21st amino acid). Pyrrolysine (considered the 22nd amino acid), which is proteinogenic only in certain microorganisms, is not used by and therefore non-essential for most organisms, including humans.
The limiting amino acid is the essential amino acid which is furthest from meeting nutritional requirements. This concept is important when determining the selection, number, and amount of foods to consume because even when total protein and all other essential amino acids are satisfied if the limiting amino acid is not satisfied then the meal is considered to be nutritionally limited by that amino acid.
Overview
(*) Pyrrolysine, sometimes considered the "22nd amino acid", is not used by the human body.
Essentiality in humans
Of the twenty amino acids common to all life forms (not counting selenocysteine), humans cannot synthesize nine: histidine, isoleucine, leucine, lysine, methionine, phenylalanine, threonine, tryptophan and valine. Additionally, the amino acids arginine, cysteine, glutamine, glycine, proline and tyrosine are considered conditionally essential, which means that specific populations who do not synthesize it in adequate amounts, such as newborn infants and people with diseased livers who are unable to synthesize cysteine, must obtain one or more of these conditionally essential amino acids from their diet. For example, enough arginine is synthesized by the urea cycle to meet the needs of an adult but perhaps not those of a growing child. Amino acids that must be obtained from the diet are called essential amino acids.
Eukaryotes can synthesize some of the amino acids from other substrates. Consequently, only a subset of the amino acids used in protein synthesis are essential nutrients.
From intermediates of the citric acid cycle and other pathways
Nonessential amino acids are produced in the body. The pathways for the synthesis of nonessential amino acids come from basic metabolic pathways. Glutamate dehydrogenase catalyzes the reductive amination of α-ketoglutarate to glutamate. A transamination reaction takes place in the synthesis of most amino acids. At this step, the chirality of the amino acid is established. Alanine and aspartate are synthesized by the transamination of pyruvate and oxaloacetate, respectively. Glutamine is synthesized from NH4+ and glutamate, and asparagine is synthesized similarly. Proline and arginine are both derived from glutamate. Serine, formed from 3-phosphoglycerate, which comes from glycolysis, is the precursor of glycine and cysteine. Tyrosine is synthesized by the hydroxylation of phenylalanine, which is an essential amino acid.
Recommended daily intake
Estimating the daily requirement for the indispensable amino acids has proven to be difficult; these numbers have undergone considerable revision over the last 20 years. The following table lists the recommended daily amounts currently in use for essential amino acids in adult humans (unless specified otherwise), together with their standard one-letter abbreviations.
The recommended daily intakes for children aged three years and older is 10% to 20% higher than adult levels and those for infants can be as much as 150% higher in the first year of life. Cysteine (or sulfur-containing amino acids), tyrosine (or aromatic amino acids), and arginine are always required by infants and growing children. Methionine and cysteine are grouped together because one of them can be synthesized from the other using the enzyme methionine S-methyltransferase and the catalyst methionine synthase. Phenylalanine and tyrosine are grouped together because tyrosine can be synthesized from phenylalanine using the enzyme phenylalanine hydroxylase.
Amino acid requirements and the amino acid content of food
Historically, amino acid requirements were determined by calculating the balance between dietary nitrogen intake and nitrogen excreted in the liquid and solid wastes, because proteins represent the largest nitrogen content in a body. A positive balance occurs when more nitrogen is consumed than is excreted, which indicates that some of the nitrogen is being used by the body to build proteins. A negative nitrogen balance occurs when more nitrogen is excreted than is consumed, which indicates that there is insufficient intake for the body to maintain its health. Graduate students at the University of Illinois were fed an artificial diet so that there was a slightly positive nitrogen balance. Then one amino acid was omitted and the nitrogen balance recorded. If a positive balance continued, then that amino acid was deemed not essential. If a negative balance occurred, then that amino acid was slowly restored until a slightly positive nitrogen balance stabilized and the minimum amount recorded.
A similar method was used to determine the protein content of foods. Test subjects were fed a diet containing no protein and the nitrogen losses recorded. During the first week or more there is a rapid loss of labile proteins. Once the nitrogen losses stabilize, this baseline is determined to be the minimum required for maintenance. Then the test subjects were fed a measured amount of the food being tested. The difference between the nitrogen in that food and the nitrogen losses above baseline was the amount the body retained to rebuild proteins. The amount of nitrogen retained divided by the total nitrogen intake is called net protein utilization. The amount of nitrogen retained divided by the (nitrogen intake minus nitrogen loss above baseline) is called biological value and is usually given as a percentage.
Modern techniques make use of ion exchange chromatography to determine the actual amino acid content of foods. The USDA used this technique in their own labs to determine the content of 7793 foods across 28 categories. The USDA published the final database in 2018 to the public.
The limiting amino acid depends on the human requirements and there are currently two sets of human requirements from authoritative sources: one published by WHO and the other published by USDA.
Protein quality
Various attempts have been made to express the "quality" or "value" of various kinds of protein. Measures include the biological value, net protein utilization, protein efficiency ratio, protein digestibility corrected amino acid score and the complete proteins concept. These concepts are important in the livestock industry, because the relative lack of one or more of the essential amino acids in animal feeds would have a limiting effect on growth and thus on feed conversion ratio. Thus, various feedstuffs may be fed in combination to increase net protein utilization, or a supplement of an individual amino acid (methionine, lysine, threonine, or tryptophan) can be added to the feed.
Protein per calorie
Protein content in foods is often measured in protein per serving rather than protein per calorie. For instance, the USDA lists 6 grams of protein per large whole egg (a 50-gram serving) rather than 84 mg of protein per calorie (71 calories total). For comparison, there are 2.8 grams of protein in a serving of raw broccoli (100 grams) or 82 mg of protein per calorie (34 calories total), or the Daily Value of 47.67g of protein after eating 1,690g of raw broccoli a day at 574 cal. An egg contains 12.5g of protein per 100g, but 4 mg more protein per calorie, or the protein DV after 381g of egg, which is 545 cal. The ratio of essential amino acids (the quality of protein) is not taken into account, one would actually need to eat more than 3 kg of broccoli a day to have a healthy protein profile, and almost 6 kg to get enough calories. It is recommended that adult humans obtain between 10–35% of their 2000 calories a day as protein.
Complete proteins in non-human animals
Scientists had known since the early 20th century that rats could not survive on a diet whose only protein source was zein, which comes from maize (corn), but recovered if they were fed casein from cow's milk. This led William Cumming Rose to the discovery of the essential amino acid threonine. Through manipulation of rodent diets, Rose was able to show that ten amino acids are essential for rats: lysine, tryptophan, histidine, phenylalanine, leucine, isoleucine, methionine, valine, and arginine, in addition to threonine. Rose's later work showed that eight amino acids are essential for adult human beings, with histidine also being essential for infants. Longer-term studies established histidine as also essential for adult humans.
Interchangeability
The distinction between essential and non-essential amino acids is somewhat unclear, as some amino acids can be produced from others. The sulfur-containing amino acids, methionine and homocysteine, can be converted into each other but neither can be synthesized de novo in humans. Likewise, cysteine can be made from homocysteine but cannot be synthesized on its own. So, for convenience, sulfur-containing amino acids are sometimes considered a single pool of nutritionally equivalent amino acids as are the aromatic amino acid pair, phenylalanine and tyrosine. Likewise arginine, ornithine, and citrulline, which are interconvertible by the urea cycle, are considered a single group.
Effects of deficiency
If one of the essential amino acids is not available in the required quantities, protein synthesis will be inhibited, irrespective of the availability of the other amino acids. Protein deficiency has been shown to affect all of the body's organs and many of its systems, for example affecting brain development in infants and young children; inhibiting upkeep of the immune system, increasing risk of infection; affecting gut mucosal function and permeability, thereby reducing absorption and increasing vulnerability to systemic disease; and impacting kidney function. The physical signs of protein deficiency include edema, failure to thrive in infants and children, poor musculature, dull skin, and thin and fragile hair. Biochemical changes reflecting protein deficiency include low serum albumin and low serum transferrin.
The amino acids that are essential in the human diet were established in a series of experiments led by William Cumming Rose. The experiments involved elemental diets to healthy male graduate students. These diets consisted of corn starch, sucrose, butterfat without protein, corn oil, inorganic salts, the known vitamins, a large brown "candy" made of liver extract flavored with peppermint oil (to supply any unknown vitamins), and mixtures of highly purified individual amino acids. The main outcome measure was nitrogen balance. Rose noted that the symptoms of nervousness, exhaustion, and dizziness were encountered to a greater or lesser extent whenever human subjects were deprived of an essential amino acid.
Essential amino acid deficiency should be distinguished from protein-energy malnutrition, which can manifest as marasmus or kwashiorkor. Kwashiorkor was once attributed to pure protein deficiency in individuals who were consuming enough calories ("sugar baby syndrome"). However, this theory has been challenged by the finding that there is no difference in the diets of children developing marasmus as opposed to kwashiorkor. Still, for instance in Dietary Reference Intakes (DRI) maintained by the USDA, lack of one or more of the essential amino acids is described as protein-energy malnutrition.
| Biology and health sciences | Amino acids | Biology |
102359 | https://en.wikipedia.org/wiki/Pulmonary%20alveolus | Pulmonary alveolus | A pulmonary alveolus (; ), also called an air sac or air space, is one of millions of hollow, distensible cup-shaped cavities in the lungs where pulmonary gas exchange takes place. Oxygen is exchanged for carbon dioxide at the blood–air barrier between the alveolar air and the pulmonary capillary. Alveoli make up the functional tissue of the mammalian lungs known as the lung parenchyma, which takes up 90 percent of the total lung volume.
Alveoli are first located in the respiratory bronchioles that mark the beginning of the respiratory zone. They are located sparsely in these bronchioles, line the walls of the alveolar ducts, and are more numerous in the blind-ended alveolar sacs. The acini are the basic units of respiration, with gas exchange taking place in all the alveoli present. The alveolar membrane is the gas exchange surface, surrounded by a network of capillaries. Oxygen is diffused across the membrane into the capillaries and carbon dioxide is released from the capillaries into the alveoli to be breathed out.
Alveoli are particular to mammalian lungs. Different structures are involved in gas exchange in other vertebrates.
Structure
The alveoli are first located in the respiratory bronchioles as scattered outpockets, extending from their lumens. The respiratory bronchioles run for considerable lengths and become increasingly alveolated with side branches of alveolar ducts that become deeply lined with alveoli. The ducts number between two and eleven from each bronchiole. Each duct opens into five or six alveolar sacs into which clusters of alveoli open.
Each terminal respiratory unit is called an acinus and consists of the respiratory bronchioles, alveolar ducts, alveolar sacs, and alveoli. New alveoli continue to form until the age of eight years.
A typical pair of human lungs contains about 480 million alveoli, providing a total surface area for gas exchange of between 70 and 80 square metres. Each alveolus is wrapped in a fine mesh of capillaries covering about 70% of its area. The diameter of an alveolus is between 200 and 500 μm.
Microanatomy
An alveolus consists of an epithelial layer of simple squamous epithelium (very thin, flattened cells), and an extracellular matrix surrounded by capillaries. The epithelial lining is part of the alveolar membrane, also known as the respiratory membrane, that allows the exchange of gases. The membrane has several layers – a layer of alveolar lining fluid that contains surfactant, the epithelial layer and its basement membrane; a thin interstitial space between the epithelial lining and the capillary membrane; a capillary basement membrane that often fuses with the alveolar basement membrane, and the capillary endothelial membrane. The whole membrane however is only between 0.2 μm at its thinnest part and 0.6 μm at its thickest.
In the alveolar walls there are interconnecting air passages between the alveoli known as the pores of Kohn. The alveolar septum that separates the alveoli in the alveolar sac contains some collagen fibers and elastic fibers. The septa also house the enmeshed capillary network that surrounds each alveolus. The elastic fibres allow the alveoli to stretch when they fill with air during inhalation. They then spring back during exhalation in order to expel the carbon dioxide-rich air.
There are three major types of alveolar cell. Two types are pneumocytes or pneumonocytes known as type I and type II cells found in the alveolar wall, and a large phagocytic cell known as an alveolar macrophage that moves about in the lumens of the alveoli, and in the connective tissue between them. Type I cells, also called type I pneumocytes, or type I alveolar cells, are squamous, thin and flat and form the structure of the alveoli. Type II cells, also called type II pneumocytes or type II alveolar cells, release pulmonary surfactant to lower surface tension, and can also differentiate to replace damaged type I cells.
Development
Development of the earliest structures that will contain alveoli begins on day 22 and is divided into five stages: embryonic, pseudoglandular, canalicular, saccular, and alveolar stage. The alveolar stage begins approximately 36 weeks into development. Immature alveoli appear as bulges from the sacculi which invade the primary septa. As the sacculi develop, the protrusions in the primary septa become larger; new septations are longer and thinner and are known as secondary septa. Secondary septa are responsible for the final division of the sacculi into alveoli. Majority of alveolar division occurs within the first 6 months but continue to develop until 3 years of age. To create a thinner diffusion barrier, the double-layer capillary network fuse into one network, each one closely associated with two alveoli as they develop.
In the first three years of life, the enlargement of lungs is a consequence of the increasing number of alveoli; after this point, both the number and size of alveoli increases until the development of lungs finishes at approximately 8 years of age.
Function
Type I cells
Type I cells are the larger of the two cell types; they are thin, flat epithelial lining cells (membranous pneumocytes), that form the structure of the alveoli. They are squamous (giving more surface area to each cell) and have long cytoplasmic extensions that cover more than 95% of the alveolar surface.
Type I cells are involved in the process of gas exchange between the alveoli and blood. These cells are extremely thin – sometimes only 25 nm – the electron microscope was needed to prove that all alveoli are lined with epithelium. This thin lining enables a fast diffusion of gas exchange between the air in the alveoli and the blood in the surrounding capillaries.
The nucleus of a type I cell occupies a large area of free cytoplasm and its organelles are clustered around it reducing the thickness of the cell. This also keeps the thickness of the blood-air barrier reduced to a minimum.
The cytoplasm in the thin portion contains pinocytotic vesicles which may play a role in the removal of small particulate contaminants from the outer surface. In addition to desmosomes, all type I alveolar cells have occluding junctions that prevent the leakage of tissue fluid into the alveolar air space.
The relatively low solubility (and hence rate of diffusion) of oxygen necessitates the large internal surface area (about 80 square m [96 square yards]) and very thin walls of the alveoli. Weaving between the capillaries and helping to support them is an extracellular matrix, a meshlike fabric of elastic and collagenous fibres. The collagen fibres, being more rigid, give the wall firmness, while the elastic fibres permit expansion and contraction of the walls during breathing.
Type I pneumocytes are unable to replicate and are susceptible to toxic insults. In the event of damage, type II cells can proliferate and differentiate into type I cells to compensate.
Type II cells
Type II cells are cuboidal and much smaller than type I cells. They are the most numerous cells in the alveoli, yet do not cover as much surface area as the squamous type I cells. Type II cells (granulous pneumocytes) in the alveolar wall contain secretory organelles known as lamellar bodies or lamellar granules, that fuse with the cell membranes and secrete pulmonary surfactant. This surfactant is a film of fatty substances, a group of phospholipids that reduce alveolar surface tension. The phospholipids are stored in the lamellar bodies. Without this coating, the alveoli would collapse. The surfactant is continuously released by exocytosis. Reinflation of the alveoli following exhalation is made easier by the surfactant, which reduces surface tension in the thin fluid lining of the alveoli. The fluid coating is produced by the body in order to facilitate the transfer of gases between blood and alveolar air, and the type II cells are typically found at the blood–air barrier.
Type II cells start to develop at about 26 weeks of gestation, secreting small amounts of surfactant. However, adequate amounts of surfactant are not secreted until about 35 weeks of gestation – this is the main reason for increased rates of infant respiratory distress syndrome, which drastically reduces at ages above 35 weeks gestation.
Type II cells are also capable of cellular division, giving rise to more type I and II alveolar cells when the lung tissue is damaged.
MUC1, a human gene associated with type II pneumocytes, has been identified as a marker in lung cancer.
The importance of the type 2 lung alveolar cells in the development of severe respiratory symptoms of COVID-19 and potential mechanisms on how these cells are protected by the SSRIs fluvoxamine and fluoxetine was summarized in a review in April 2022.
Alveolar macrophages
The alveolar macrophages reside on the internal luminal surfaces of the alveoli, the alveolar ducts, and the bronchioles. They are mobile scavengers that serve to engulf foreign particles in the lungs, such as dust, bacteria, carbon particles, and blood cells from injuries. They are also called pulmonary macrophages, and dust cells. Alveolar macrophages also play a crucial role in immune responses against viral pathogens in the lungs. They secrete cytokines and chemokines, which recruit and activate other immune cells, initiate type I interferon signaling, and inhibit the nuclear export of viral genomes.
Clinical significance
Diseases
Surfactant
Insufficient surfactant in the alveoli is one of the causes that can contribute to atelectasis (collapse of part or all of the lung). Without pulmonary surfactant, atelectasis is a certainty. The severe condition of acute respiratory distress syndrome (ARDS) is caused by a deficiency or dysfunction of surfactant. Insufficient surfactant in the lungs of preterm infants causes infant respiratory distress syndrome (IRDS). The lecithin–sphingomyelin ratio is a measure of fetal amniotic fluid to indicate lung maturity or immaturity. A low ratio indicates a risk factor for IRDS. Lecithin and sphingomyelin are two of the glycolipids of pulmonary surfactant.
Impaired surfactant regulation can cause an accumulation of surfactant proteins to build up in the alveoli in a condition called pulmonary alveolar proteinosis. This results in impaired gas exchange.
Inflammation
Pneumonia is an inflammatory condition of the lung tissue, which can be caused by both viruses and bacteria. Cytokines and fluids are released into the alveolar cavity, interstitium, or both, in response to infection, causing the effective surface area of gas exchange to be reduced. In severe cases where cellular respiration cannot be maintained, supplemental oxygen may be required.
Diffuse alveolar damage can be a cause of acute respiratory distress syndrome(ARDS) a severe inflammatory disease of the lung.
In asthma, the bronchioles become narrowed, causing the amount of air flow into the lung tissue to be greatly reduced. It can be triggered by irritants in the air, photochemical smog for example, as well as substances to which a person is allergic.
Chronic bronchitis occurs when an abundance of mucus is produced by the lungs. The production of mucus occurs naturally when the lung tissue is exposed to irritants. In chronic bronchitis, the air passages into the alveoli, the respiratory bronchioles, become clogged with mucus. This causes increased coughing in order to remove the mucus, and is often a result of extended periods of exposure to cigarette smoke.
Hypersensitivity pneumonitis
Structural
Almost any type of lung tumor or lung cancer can compress the alveoli and reduce gas exchange capacity. In some cases the tumor will fill the alveoli.
Cavitary pneumonia is a process in which the alveoli are destroyed and produce a cavity. As the alveoli are destroyed, the surface area for gas exchange to occur becomes reduced. Further changes in blood flow can lead to decline in lung function.
Emphysema is another disease of the lungs, whereby the elastin in the walls of the alveoli is broken down by an imbalance between the production of neutrophil elastase (elevated by cigarette smoke) and alpha-1 antitrypsin (the activity varies due to genetics or reaction of a critical methionine residue with toxins including cigarette smoke). The resulting loss of elasticity in the lungs leads to prolonged times for exhalation, which occurs through passive recoil of the expanded lung. This leads to a smaller volume of gas exchanged per breath.
Pulmonary alveolar microlithiasis is a rare lung disorder of small stone formation in the alveoli.
Several factors, including smoking, viral infections, and aging, contribute to physical damage to type II alveolar cells. Some studies have linked injury to these cells to the proliferation of fibrosis in the lungs and the onset of idiopathic pulmonary fibrosis.
Fluid
A pulmonary contusion is a bruise of the lung tissue caused by trauma. Damaged capillaries from a contusion can cause blood and other fluids to accumulate in the tissue of the lung, impairing gas exchange.
Pulmonary edema is the buildup of fluid in the parenchyma and alveoli. An edema is usually caused by left ventricular heart failure, or by damage to the lung or its vasculature.
Coronavirus
Because of the high expression of angiotensin-converting enzyme 2 (ACE2) in type II alveolar cells, the lungs are susceptible to infections by some coronaviruses including the viruses that cause severe acute respiratory syndrome (SARS) and coronavirus disease 2019 (COVID-19).
Additional images
| Biology and health sciences | Respiratory system | Biology |
102519 | https://en.wikipedia.org/wiki/Dolomite%20%28mineral%29 | Dolomite (mineral) | Dolomite () is an anhydrous carbonate mineral composed of calcium magnesium carbonate, ideally The term is also used for a sedimentary carbonate rock composed mostly of the mineral dolomite (see Dolomite (rock)). An alternative name sometimes used for the dolomitic rock type is dolostone.
History
As stated by Nicolas-Théodore de Saussure the mineral dolomite was probably first described by Carl Linnaeus in 1768. In 1791, it was described as a rock by the French naturalist and geologist Déodat Gratet de Dolomieu (1750–1801), first in buildings of the old city of Rome, and later as samples collected in the Tyrolean Alps. Nicolas-Théodore de Saussure first named the mineral (after Dolomieu) in March 1792.
Properties
The mineral dolomite crystallizes in the trigonal-rhombohedral system. It forms white, tan, gray, or pink crystals. Dolomite is a double carbonate, having an alternating structural arrangement of calcium and magnesium ions. Unless it is in fine powder form, it does not rapidly dissolve or effervesce (fizz) in cold dilute hydrochloric acid as calcite does. Crystal twinning is common.
Solid solution exists between dolomite, the iron-dominant ankerite and the manganese-dominant kutnohorite. Small amounts of iron in the structure give the crystals a yellow to brown tint. Manganese substitutes in the structure also up to about three percent MnO. A high manganese content gives the crystals a rosy pink color. Lead, zinc, and cobalt also can substitute in the structure for magnesium. The mineral dolomite is closely related to huntite .
Because dolomite can be dissolved by slightly acidic water, areas where dolomite is an abundant rock-forming mineral are important as aquifers and contribute to karst terrain formation.
Formation
Modern dolomite formation has been found to occur under anaerobic conditions in supersaturated saline lagoons such as those at the Rio de Janeiro coast of Brazil, namely, Lagoa Vermelha and Brejo do Espinho. There are many other localities where modern dolomite forms, notably along sabkhas in the Persian Gulf, but also in sedimentary basins bearing gas hydrates and hypersaline lakes. It is often thought that dolomite nucleates with the help of sulfate-reducing bacteria (e.g. Desulfovibrio brasiliensis), but other microbial metabolisms have been also found to mediate in dolomite formation. In general, low-temperature dolomite may occur in natural supersaturated environments rich in extracellular polymeric substances (EPS) and microbial cell surfaces. This is likely result from complexation of both magnesium and calcium by carboxylic acids comprising EPS.
Vast deposits of dolomite are present in the geological record, but the mineral is relatively rare in modern environments. Reproducible, inorganic low-temperature syntheses of dolomite are yet to be performed. Usually, the initial inorganic precipitation of a metastable "precursor" (such as magnesium calcite) can easily be achieved. The precursor phase will theoretically change gradually into a more stable phase (such as partially ordered dolomite) during periodical intervals of dissolution and re-precipitation. The general principle governing the course of this irreversible geochemical reaction has been coined "breaking Ostwald's step rule". High diagenetic temperatures, such as those of groundwater flowing along deeply rooted fault systems affecting some sedimentary successions or deeply buried limestone rocks allocate dolomitization. But the mineral is also volumetrically important in some Neogene platforms never subjected to elevated temperatures. Under such conditions of diagenesis the long-term activity of the deep biosphere could play a key role in dolomitization, since diagenetic fluids of contrasting composition are mixed as a response to Milankovitch cycles.
A recent biotic synthetic experiment claims to have precipitated ordered dolomite when anoxygenic photosynthesis proceeds in the presence of manganese(II). A still perplexing example of an organogenic origin is that of the reported formation of dolomite in the urinary bladder of a Dalmatian dog, possibly as the result of an illness or infection.
Uses
Dolomite is used as an ornamental stone, a concrete aggregate, and a source of magnesium oxide, as well as in the Pidgeon process for the production of magnesium. It is an important petroleum reservoir rock, and serves as the host rock for large strata-bound Mississippi Valley-Type (MVT) ore deposits of base metals such as lead, zinc, and copper. Where calcite limestone is uncommon or too costly, dolomite is sometimes used in its place as a flux for the smelting of iron and steel. Large quantities of processed dolomite are used in the production of float glass.
In horticulture, dolomite and dolomitic limestone are added to soils and soilless potting mixes as a pH buffer and as a magnesium source. Pastures can be limed with dolomitic lime to raise their pH and where there is a magnesium deficiency.
Dolomite is also used as the substrate in marine (saltwater) aquariums to help buffer changes in the pH of the water.
Calcined dolomite is also used as a catalyst for destruction of tar in the gasification of biomass at high temperature. Particle physics researchers like to build particle detectors under layers of dolomite to enable the detectors to detect the highest possible number of exotic particles. Because dolomite contains relatively minor quantities of radioactive materials, it can insulate against interference from cosmic rays without adding to background radiation levels.
In addition to being an industrial mineral, dolomite is highly valued by collectors and museums when it forms large, transparent crystals. The specimens that appear in the magnesite quarry exploited in Eugui, Esteribar, Navarra (Spain) are considered among the best in the world.
| Physical sciences | Minerals | Earth science |
10228498 | https://en.wikipedia.org/wiki/Lissajous%20orbit | Lissajous orbit | In orbital mechanics, a Lissajous orbit (), named after Jules Antoine Lissajous, is a quasi-periodic orbital trajectory that an object can follow around a Lagrangian point of a three-body system with minimal propulsion. Lyapunov orbits around a Lagrangian point are curved paths that lie entirely in the plane of the two primary bodies. In contrast, Lissajous orbits include components in this plane and perpendicular to it, and follow a Lissajous curve. Halo orbits also include components perpendicular to the plane, but they are periodic, while Lissajous orbits are usually not.
In practice, any orbits around Lagrangian points , , or are dynamically unstable, meaning small departures from equilibrium grow over time. As a result, spacecraft in these Lagrangian point orbits must use their propulsion systems to perform orbital station-keeping. Although they are not perfectly stable, a modest effort of station keeping keeps a spacecraft in a desired Lissajous orbit for a long time.
In the absence of other influences, orbits about Lagrangian points and are dynamically stable so long as the ratio of the masses of the two main objects is greater than about 25. The natural dynamics keep the spacecraft (or natural celestial body) in the vicinity of the Lagrangian point without use of a propulsion system, even when slightly perturbed from equilibrium. These orbits can however be destabilized by other nearby massive objects. For example, orbits around the and points in the Earth–Moon system can last only a few million years instead of billions because of perturbations by the other planets in the Solar System.
Spacecraft using Lissajous orbits
Several missions have used Lissajous orbits: ACE at Sun–Earth L1, SOHO at Sun–Earth L1, DSCOVR at Sun–Earth L1, WMAP at Sun–Earth L2, and also the Genesis mission collecting solar particles at L1.
On 14 May 2009, the European Space Agency (ESA) launched into space the Herschel and Planck observatories, both of which use Lissajous orbits at Sun–Earth L2.
ESA's Gaia mission also uses a Lissajous orbit at Sun–Earth L2.
In 2011, NASA transferred two of its THEMIS spacecraft from Earth orbit to Lunar orbit by way of Earth–Moon L1 and L2 Lissajous orbits.
In June 2018, Queqiao, the relay satellite for China's Chang'e 4 lunar lander mission, entered orbit around Earth-Moon L2.
Fictional appearances
In the 2005 science fiction novel Sunstorm by Arthur C. Clarke and Stephen Baxter, a huge shield is constructed in space to protect the Earth from a deadly solar storm. The shield is described to have been in a Lissajous orbit at . In the story a group of wealthy and powerful people shelter opposite the shield at so as to be protected from the solar storm by the shield, the Earth and the Moon.
In the 2017 science fiction novel Artemis by Andy Weir, a Lissajous orbit is used as a transfer point for routine travel to and from the Moon.
| Physical sciences | Orbital mechanics | Astronomy |
15368428 | https://en.wikipedia.org/wiki/Radio | Radio | Radio is the technology of communicating using radio waves. Radio waves are electromagnetic waves of frequency between 3 hertz (Hz) and 300 gigahertz (GHz). They are generated by an electronic device called a transmitter connected to an antenna which radiates oscillating electrical energy, often characterized as a wave. They can be received by other antennas connected to a radio receiver; this is the fundamental principle of radio communication. In addition to communication, radio is used for radar, radio navigation, remote control, remote sensing, and other applications.
In radio communication, used in radio and television broadcasting, cell phones, two-way radios, wireless networking, and satellite communication, among numerous other uses, radio waves are used to carry information across space from a transmitter to a receiver, by modulating the radio signal (impressing an information signal on the radio wave by varying some aspect of the wave) in the transmitter. In radar, used to locate and track objects like aircraft, ships, spacecraft and missiles, a beam of radio waves emitted by a radar transmitter reflects off the target object, and the reflected waves reveal the object's location to a receiver that is typically colocated with the transmitter. In radio navigation systems such as GPS and VOR, a mobile navigation instrument receives radio signals from multiple navigational radio beacons whose position is known, and by precisely measuring the arrival time of the radio waves the receiver can calculate its position on Earth. In wireless radio remote control devices like drones, garage door openers, and keyless entry systems, radio signals transmitted from a controller device control the actions of a remote device.
The existence of radio waves was first proven by German physicist Heinrich Hertz on 11 November 1886. In the mid-1890s, building on techniques physicists were using to study electromagnetic waves, Italian physicist Guglielmo Marconi developed the first apparatus for long-distance radio communication, sending a wireless Morse Code message to a recipient over a kilometer away in 1895, and the first transatlantic signal on 12 December 1901. The first commercial radio broadcast was transmitted on 2 November 1920, when the live returns of the Harding-Cox presidential election were broadcast by Westinghouse Electric and Manufacturing Company in Pittsburgh, under the call sign KDKA.
The emission of radio waves is regulated by law, coordinated by the International Telecommunication Union (ITU), which allocates frequency bands in the radio spectrum for various uses.
Etymology
The word radio is derived from the Latin word radius, meaning "spoke of a wheel, beam of light, ray". It was first applied to communications in 1881 when, at the suggestion of French scientist , Alexander Graham Bell adopted radiophone (meaning "radiated sound") as an alternate name for his photophone optical transmission system.
Following Hertz's discovery of the existence of radio waves in 1886, the term Hertzian waves was initially used for this radiation. The first practical radio communication systems, developed by Marconi in 1894–1895, transmitted telegraph signals by radio waves, so radio communication was first called wireless telegraphy. Up until about 1910 the term wireless telegraphy also included a variety of other experimental systems for transmitting telegraph signals without wires, including electrostatic induction, electromagnetic induction and aquatic and earth conduction, so there was a need for a more precise term referring exclusively to electromagnetic radiation.
The French physicist Édouard Branly, who in 1890 developed the radio wave detecting coherer, called it in French a radio-conducteur. The radio- prefix was later used to form additional descriptive compound and hyphenated words, especially in Europe. For example, in early 1898 the British publication The Practical Engineer included a reference to the radiotelegraph and radiotelegraphy.
The use of radio as a standalone word dates back to at least 30 December 1904, when instructions issued by the British Post Office for transmitting telegrams specified that "The word 'Radio'... is sent in the Service Instructions." This practice was universally adopted, and the word "radio" introduced internationally, by the 1906 Berlin Radiotelegraphic Convention, which included a Service Regulation specifying that "Radiotelegrams shall show in the preamble that the service is 'Radio.
The switch to radio in place of wireless took place slowly and unevenly in the English-speaking world. Lee de Forest helped popularize the new word in the United States—in early 1907, he founded the DeForest Radio Telephone Company, and his letter in the 22 June 1907 Electrical World about the need for legal restrictions warned that "Radio chaos will certainly be the result until such stringent regulation is enforced." The United States Navy would also play a role. Although its translation of the 1906 Berlin Convention used the terms wireless telegraph and wireless telegram, by 1912 it began to promote the use of radio instead. The term started to become preferred by the general public in the 1920s with the introduction of broadcasting.
History
Electromagnetic waves were predicted by James Clerk Maxwell in his 1873 theory of electromagnetism, now called Maxwell's equations, who proposed that a coupled oscillating electric field and magnetic field could travel through space as a wave, and proposed that light consisted of electromagnetic waves of short wavelength. On 11 November 1886, German physicist Heinrich Hertz, attempting to confirm Maxwell's theory, first observed radio waves he generated using a primitive spark-gap transmitter. Experiments by Hertz and physicists Jagadish Chandra Bose, Oliver Lodge, Lord Rayleigh, and Augusto Righi, among others, showed that radio waves like light demonstrated reflection, refraction, diffraction, polarization, standing waves, and traveled at the same speed as light, confirming that both light and radio waves were electromagnetic waves, differing only in frequency. In 1895, Guglielmo Marconi developed the first radio communication system, using a spark-gap transmitter to send Morse code over long distances. By December 1901, he had transmitted across the Atlantic Ocean. Marconi and Karl Ferdinand Braun shared the 1909 Nobel Prize in Physics "for their contributions to the development of wireless telegraphy".
During radio's first two decades, called the radiotelegraphy era, the primitive radio transmitters could only transmit pulses of radio waves, not the continuous waves which were needed for audio modulation, so radio was used for person-to-person commercial, diplomatic and military text messaging. Starting around 1908 industrial countries built worldwide networks of powerful transoceanic transmitters to exchange telegram traffic between continents and communicate with their colonies and naval fleets. During World War I the development of continuous wave radio transmitters, rectifying electrolytic, and crystal radio receiver detectors enabled amplitude modulation (AM) radiotelephony to be achieved by Reginald Fessenden and others, allowing audio to be transmitted. On 2 November 1920, the first commercial radio broadcast was transmitted by Westinghouse Electric and Manufacturing Company in Pittsburgh, under the call sign KDKA featuring live coverage of the Harding-Cox presidential election.
Technology
Radio waves are radiated by electric charges undergoing acceleration. They are generated artificially by time-varying electric currents, consisting of electrons flowing back and forth in a metal conductor called an antenna.
As they travel farther from the transmitting antenna, radio waves spread out so their signal strength (intensity in watts per square meter) decreases (see Inverse-square law), so radio transmissions can only be received within a limited range of the transmitter, the distance depending on the transmitter power, the antenna radiation pattern, receiver sensitivity, background noise level, and presence of obstructions between transmitter and receiver. An omnidirectional antenna transmits or receives radio waves in all directions, while a directional antenna transmits radio waves in a beam in a particular direction, or receives waves from only one direction.
Radio waves travel at the speed of light in vacuum and at slightly lower velocity in air.
The other types of electromagnetic waves besides radio waves, infrared, visible light, ultraviolet, X-rays and gamma rays, can also carry information and be used for communication. The wide use of radio waves for telecommunication is mainly due to their desirable propagation properties stemming from their longer wavelength.
Radio communication
In radio communication systems, information is carried across space using radio waves. At the sending end, the information to be sent is converted by some type of transducer to a time-varying electrical signal called the modulation signal. The modulation signal may be an audio signal representing sound from a microphone, a video signal representing moving images from a video camera, or a digital signal consisting of a sequence of bits representing binary data from a computer. The modulation signal is applied to a radio transmitter. In the transmitter, an electronic oscillator generates an alternating current oscillating at a radio frequency, called the carrier wave because it serves to generate the radio waves that carry the information through the air. The modulation signal is used to modulate the carrier, varying some aspect of the carrier wave, impressing the information in the modulation signal onto the carrier. Different radio systems use different modulation methods:
Amplitude modulation (AM) – in an AM transmitter, the amplitude (strength) of the radio carrier wave is varied by the modulation signal;
Frequency modulation (FM) – in an FM transmitter, the frequency of the radio carrier wave is varied by the modulation signal;
Frequency-shift keying (FSK) – used in wireless digital devices to transmit digital signals, the frequency of the carrier wave is shifted between frequencies.
orthogonal frequency-division multiplexing (OFDM) – a family of digital modulation methods widely used in high-bandwidth systems such as Wi-Fi networks, cellphones, digital television broadcasting, and digital audio broadcasting (DAB) to transmit digital data using a minimum of radio spectrum bandwidth. It has higher spectral efficiency and more resistance to fading than AM or FM. In OFDM, multiple radio carrier waves closely spaced in frequency are transmitted within the radio channel, with each carrier modulated with bits from the incoming bitstream so multiple bits are being sent simultaneously, in parallel. At the receiver, the carriers are demodulated and the bits are combined in the proper order into one bitstream.
Many other types of modulation are also used. In some types, the carrier wave is suppressed, and only one or both modulation sidebands are transmitted.
The modulated carrier is amplified in the transmitter and applied to a transmitting antenna which radiates the energy as radio waves. The radio waves carry the information to the receiver location. At the receiver, the radio wave induces a tiny oscillating voltage in the receiving antennaa weaker replica of the current in the transmitting antenna. This voltage is applied to the radio receiver, which amplifies the weak radio signal so it is stronger, then demodulates it, extracting the original modulation signal from the modulated carrier wave. The modulation signal is converted by a transducer back to a human-usable form: an audio signal is converted to sound waves by a loudspeaker or earphones, a video signal is converted to images by a display, while a digital signal is applied to a computer or microprocessor, which interacts with human users.
The radio waves from many transmitters pass through the air simultaneously without interfering with each other because each transmitter's radio waves oscillate at a different rate, in other words, each transmitter has a different frequency, measured in hertz (Hz), kilohertz (kHz), megahertz (MHz) or gigahertz (GHz). The receiving antenna typically picks up the radio signals of many transmitters. The receiver uses tuned circuits to select the radio signal desired out of all the signals picked up by the antenna and reject the others. A tuned circuit (also called resonant circuit or tank circuit) acts like a resonator, similar to a tuning fork. It has a natural resonant frequency at which it oscillates. The resonant frequency of the receiver's tuned circuit is adjusted by the user to the frequency of the desired radio station; this is called "tuning". The oscillating radio signal from the desired station causes the tuned circuit to resonate, oscillate in sympathy, and it passes the signal on to the rest of the receiver. Radio signals at other frequencies are blocked by the tuned circuit and not passed on.
Bandwidth
A modulated radio wave, carrying an information signal, occupies a range of frequencies. The information (modulation) in a radio signal is usually concentrated in narrow frequency bands called sidebands (SB) just above and below the carrier frequency. The width in hertz of the frequency range that the radio signal occupies, the highest frequency minus the lowest frequency, is called its bandwidth (BW). For any given signal-to-noise ratio, an amount of bandwidth can carry the same amount of information (data rate in bits per second) regardless of where in the radio frequency spectrum it is located, so bandwidth is a measure of information-carrying capacity. The bandwidth required by a radio transmission depends on the data rate of the information (modulation signal) being sent, and the spectral efficiency of the modulation method used; how much data it can transmit in each kilohertz of bandwidth. Different types of information signals carried by radio have different data rates. For example, a television (video) signal has a greater data rate than an audio signal.
The radio spectrum, the total range of radio frequencies that can be used for communication in a given area, is a limited resource. Each radio transmission occupies a portion of the total bandwidth available. Radio bandwidth is regarded as an economic good which has a monetary cost and is in increasing demand. In some parts of the radio spectrum, the right to use a frequency band or even a single radio channel is bought and sold for millions of dollars. So there is an incentive to employ technology to minimize the bandwidth used by radio services.
A slow transition from analog to digital radio transmission technologies began in the late 1990s. Part of the reason for this is that digital modulation can often transmit more information (a greater data rate) in a given bandwidth than analog modulation, by using data compression algorithms, which reduce redundancy in the data to be sent, and more efficient modulation. Other reasons for the transition is that digital modulation has greater noise immunity than analog, digital signal processing chips have more power and flexibility than analog circuits, and a wide variety of types of information can be transmitted using the same digital modulation.
Because it is a fixed resource which is in demand by an increasing number of users, the radio spectrum has become increasingly congested in recent decades, and the need to use it more effectively is driving many additional radio innovations such as trunked radio systems, spread spectrum (ultra-wideband) transmission, frequency reuse, dynamic spectrum management, frequency pooling, and cognitive radio.
ITU frequency bands
The ITU arbitrarily divides the radio spectrum into 12 bands, each beginning at a wavelength which is a power of ten (10n) metres, with corresponding frequency of 3 times a power of ten, and each covering a decade of frequency or wavelength. Each of these bands has a traditional name:
{| class="wikitable"
|- style="text-align:left;"
! Band name !! Abbreviation !! Frequency !! Wavelength
|- style="text-align:center;"
| Extremelylow frequency || ELF || 3–30 Hz || 100,000–10,000 km
|- style="text-align:center;"
| Superlow frequency || SLF || 30–300 Hz || 10,000 –1,000 km
|- style="text-align:center;"
| Ultralow frequency || ULF || 300–3,000 Hz || 1,000–100 km
|- style="text-align:center;"
| Verylow frequency || VLF || 3–30 kHz || 100–10 km
|- style="text-align:center;"
| Lowfrequency || LF || 30–300 kHz || 10–1 km
|- style="text-align:center;"
| Mediumfrequency || MF || 300–3,000 kHz || 1,000–100 m
|}
{| class="wikitable"
|- style="text-align:left;"
! Band name !! Abbreviation !! Frequency !! Wavelength
|- style="text-align:center;"
| Highfrequency || HF || 3–30 MHz || 100–10 m
|- style="text-align:center;"
| Veryhigh frequency || VHF || 30–300 MHz || 10–1 m
|- style="text-align:center;"
| Ultrahigh frequency || UHF || 300–3,000 MHz || 100–10 cm
|- style="text-align:center;"
| Superhigh frequency || SHF || 3–30 GHz || 10–1 cm
|- style="text-align:center;"
| Extremelyhigh frequency || EHF || 30–300 GHz || 10–1 mm
|- style="text-align:center;"
| Tremendouslyhigh frequency || THF || 300–3,000 GHz(0.3–3.0 THz) || 1.0–0.1 mm
|}
It can be seen that the bandwidth, the range of frequencies, contained in each band is not equal but increases exponentially as the frequency increases; each band contains ten times the bandwidth of the preceding band.
The term "tremendously low frequency" (TLF) has been used for wavelengths from 1–3 Hz (300,000–100,000 km), though the term has not been defined by the ITU.
Regulation
The airwaves are a resource shared by many users. Two radio transmitters in the same area that attempt to transmit on the same frequency will interfere with each other, causing garbled reception, so neither transmission may be received clearly. Interference with radio transmissions can not only have a large economic cost, but it can also be life-threatening (for example, in the case of interference with emergency communications or air traffic control).
To prevent interference between different users, the emission of radio waves is strictly regulated by national laws, coordinated by an international body, the International Telecommunication Union (ITU), which allocates bands in the radio spectrum for different uses. Radio transmitters must be licensed by governments, under a variety of license classes depending on use, and are restricted to certain frequencies and power levels. In some classes, such as radio and television broadcasting stations, the transmitter is given a unique identifier consisting of a string of letters and numbers called a call sign, which must be used in all transmissions. In order to adjust, maintain, or internally repair radiotelephone transmitters, individuals must hold a government license, such as the general radiotelephone operator license in the US, obtained by taking a test demonstrating adequate technical and legal knowledge of safe radio operation.
Exceptions to the above rules allow the unlicensed operation by the public of low power short-range transmitters in consumer products such as cell phones, cordless phones, wireless devices, walkie-talkies, citizens band radios, wireless microphones, garage door openers, and baby monitors. In the US, these fall under Part 15 of the Federal Communications Commission (FCC) regulations. Many of these devices use the ISM bands, a series of frequency bands throughout the radio spectrum reserved for unlicensed use. Although they can be operated without a license, like all radio equipment these devices generally must be type-approved before the sale.
Applications
Below are some of the most important uses of radio, organized by function.
Broadcasting
Broadcasting is the one-way transmission of information from a transmitter to receivers belonging to a public audience. Since the radio waves become weaker with distance, a broadcasting station can only be received within a limited distance of its transmitter. Systems that broadcast from satellites can generally be received over an entire country or continent. Older terrestrial radio and television are paid for by commercial advertising or governments. In subscription systems like satellite television and satellite radio the customer pays a monthly fee. In these systems, the radio signal is encrypted and can only be decrypted by the receiver, which is controlled by the company and can be deactivated if the customer does not pay.
Broadcasting uses several parts of the radio spectrum, depending on the type of signals transmitted and the desired target audience. Longwave and medium wave signals can give reliable coverage of areas several hundred kilometers across, but have a more limited information-carrying capacity and so work best with audio signals (speech and music), and the sound quality can be degraded by radio noise from natural and artificial sources. The shortwave bands have a greater potential range but are more subject to interference by distant stations and varying atmospheric conditions that affect reception.
In the very high frequency band, greater than 30 megahertz, the Earth's atmosphere has less of an effect on the range of signals, and line-of-sight propagation becomes the principal mode. These higher frequencies permit the great bandwidth required for television broadcasting. Since natural and artificial noise sources are less present at these frequencies, high-quality audio transmission is possible, using frequency modulation.
Audio: Radio broadcasting
Radio broadcasting means transmission of audio (sound) to radio receivers belonging to a public audience. Analog audio is the earliest form of radio broadcast. AM broadcasting began around 1920. FM broadcasting was introduced in the late 1930s with improved fidelity. A broadcast radio receiver is called a radio. Most radios can receive both AM and FM.
AM (amplitude modulation) – in AM, the amplitude (strength) of the radio carrier wave is varied by the audio signal. AM broadcasting, the oldest broadcasting technology, is allowed in the AM broadcast bands, between 148 and 283 kHz in the low frequency (LF) band for longwave broadcasts and between 526 and 1706 kHz in the medium frequency (MF) band for medium-wave broadcasts. Because waves in these bands travel as ground waves following the terrain, AM radio stations can be received beyond the horizon at hundreds of miles distance, but AM has lower fidelity than FM. Radiated power (ERP) of AM stations in the US is usually limited to a maximum of 10 kW, although a few (clear-channel stations) are allowed to transmit at 50 kW. AM stations broadcast in monaural audio; AM stereo broadcast standards exist in most countries, but the radio industry has failed to upgrade to them, due to lack of demand.
Shortwave broadcasting – AM broadcasting is also allowed in the shortwave bands by legacy radio stations. Since radio waves in these bands can travel intercontinental distances by reflecting off the ionosphere using skywave or "skip" propagation, shortwave is used by international stations, broadcasting to other countries.
FM (frequency modulation) – in FM the frequency of the radio carrier signal is varied slightly by the audio signal. FM broadcasting is permitted in the FM broadcast bands between about 65 and 108 MHz in the very high frequency (VHF) range. Radio waves in this band travel by line-of-sight so FM reception is limited by the visual horizon to about , and can be blocked by hills. However it is less susceptible to interference from radio noise (RFI, sferics, static), and has higher fidelity, better frequency response, and less audio distortion than AM. In the US, radiated power (ERP) of FM stations varies from 6–100 kW.
Digital radio involves a variety of standards and technologies for broadcasting digital radio signals over the air. Some systems, such as HD Radio and DRM, operate in the same wavebands as analog broadcasts, either as a replacement for analog stations or as a complementary service. Others, such as DAB/DAB+ and ISDB_Tsb, operate in wavebands traditionally used for television or satellite services.
Digital Audio Broadcasting (DAB) debuted in some countries in 1998. It transmits audio as a digital signal rather than an analog signal as AM and FM do. DAB has the potential to provide higher quality sound than FM (although many stations do not choose to transmit at such high quality), has greater immunity to radio noise and interference, makes better use of scarce radio spectrum bandwidth and provides advanced user features such as electronic program guides. Its disadvantage is that it is incompatible with previous radios so that a new DAB receiver must be purchased. Several nations have set dates to switch off analog FM networks in favor of DAB / DAB+, notably Norway in 2017 and Switzerland in 2024.
A single DAB station transmits a bandwidth signal that carries from 9–12 channels of digital audio modulated by OFDM from which the listener can choose. Broadcasters can transmit a channel at a range of different bit rates, so different channels can have different audio quality. In different countries DAB stations broadcast in either Band III (174–240 MHz) or L band (1.452–1.492 GHz) in the UHF range, so like FM reception is limited by the visual horizon to about .
HD Radio is an alternative digital radio standard widely implemented in North America. An in-band on-channel technology, HD Radio broadcasts a digital signal in a subcarrier of a station's analog FM or AM signal. Stations are able to multicast more than one audio signal in the subcarrier, supporting the transmission of multiple audio services at varying bitrates. The digital signal is transmitted using OFDM with the HDC (High-Definition Coding) proprietary audio compression format. HDC is based on, but not compatible with, the MPEG-4 standard HE-AAC. It uses a modified discrete cosine transform (MDCT) audio data compression algorithm.
Digital Radio Mondiale (DRM) is a competing digital terrestrial radio standard developed mainly by broadcasters as a higher spectral efficiency replacement for legacy AM and FM broadcasting. Mondiale means "worldwide" in French and Italian; DRM was developed in 2001, and is currently supported by 23 countries, and adopted by some European and Eastern broadcasters beginning in 2003. The DRM30 mode uses the commercial broadcast bands below 30 MHz, and is intended as a replacement for standard AM broadcast on the longwave, mediumwave, and shortwave bands. The DRM+ mode uses VHF frequencies centered around the FM broadcast band, and is intended as a replacement for FM broadcasting. It is incompatible with existing radio receivers, so it requires listeners to purchase a new DRM receiver. The modulation used is a form of OFDM called COFDM in which, up to 4 carriers are transmitted on a channel formerly occupied by a single AM or FM signal, modulated by quadrature amplitude modulation (QAM).
The DRM system is designed to be as compatible as possible with existing AM and FM radio transmitters, so that much of the equipment in existing radio stations can continue in use, augmented with DRM modulation equipment.
Satellite radio is a subscription radio service that broadcasts CD quality digital audio direct to subscribers' receivers using a microwave downlink signal from a direct broadcast communication satellite in geostationary orbit above the Earth. It is mostly intended for radios in vehicles. Satellite radio uses the 2.3 GHz S band in North America, in other parts of the world, it uses the 1.4 GHz allocated for DAB.
Audio/video: Television broadcasting
Television broadcasting is the transmission of moving images along with a synchronized audio (sound) channel by radio. The sequence of still images is displayed on a screen on a television receiver (a "television" or TV), which includes a loudspeak. Television (video) signals occupy a wider bandwidth than broadcast radio (audio) signals. Analog television, the original television technology, required 6 MHz, so the television frequency bands are divided into 6 MHz channels, now called "RF channels".
The current television standard, introduced beginning in 2006, is a digital format called high-definition television (HDTV), which transmits pictures at higher resolution, typically 1080 pixels high by 1920 pixels wide, at a rate of 25 or 30 frames per second. Digital television (DTV) transmission systems, which replaced older analog television in a transition beginning in 2006, use image compression and high-efficiency digital modulation such as OFDM and 8VSB to transmit HDTV video within a smaller bandwidth than the old analog channels, saving scarce radio spectrum space. Therefore, each of the 6 MHz analog RF channels now carries up to 7 DTV channels – these are called "virtual channels". Digital television receivers have different behavior in the presence of poor reception or noise than analog television, called the "digital cliff" effect. Unlike analog television, in which increasingly poor reception causes the picture quality to gradually degrade, in digital television picture quality is not affected by poor reception until, at a certain point, the receiver stops working and the screen goes black.
Terrestrial television, over-the-air (OTA) television, or broadcast television – the oldest television technology, is the transmission of television signals from land-based television stations to television receivers (called televisions or TVs) in viewer's homes. Terrestrial television broadcasting uses the bands 41 – 88 MHz (VHF low band or Band I, carrying RF channels 1–6), 174 – 240 MHz, (VHF high band or Band III; carrying RF channels 7–13), and 470 – 614 MHz (UHF Band IV and Band V; carrying RF channels 14 and up). The exact frequency boundaries vary in different countries. Propagation is by line-of-sight, so reception is limited by the visual horizon. In the US, the effective radiated power (ERP) of television transmitters is regulated according to height above average terrain. Viewers closer to the television transmitter can use a simple "rabbit ears" dipole antenna on top of the TV, but viewers in fringe reception areas typically require an outdoor antenna mounted on the roof to get adequate reception.
Satellite television – a set-top box which receives subscription direct-broadcast satellite television, and displays it on an ordinary television. A direct broadcast satellite in geostationary orbit above the Earth's equator transmits many channels (up to 900) modulated on a 12.2 to 12.7 GHz Ku band microwave downlink signal to a rooftop satellite dish antenna on the subscriber's residence. The microwave signal is converted to a lower intermediate frequency at the dish and conducted into the building by a coaxial cable to a set-top box connected to the subscriber's TV, where it is demodulated and displayed. The subscriber pays a monthly fee.
Time and frequency
Government standard frequency and time signal services operate time radio stations which continuously broadcast extremely accurate time signals produced by atomic clocks, as a reference to synchronize other clocks. Examples are BPC, DCF77, JJY, MSF, RTZ, TDF, WWV, and YVTO. One use is in radio clocks and watches, which include an automated receiver that periodically (usually weekly) receives and decodes the time signal and resets the watch's internal quartz clock to the correct time, thus allowing a small watch or desk clock to have the same accuracy as an atomic clock. Government time stations are declining in number because GPS satellites and the Internet Network Time Protocol (NTP) provide equally accurate time standards.
Voice communication
Two-way voice communication
A two-way radio is an audio transceiver, a receiver and transmitter in the same device, used for bidirectional person-to-person voice communication with other users with similar radios. An older term for this mode of communication is radiotelephony. The radio link may be half-duplex, as in a walkie-talkie, using a single radio channel in which only one radio can transmit at a time, so different users take turns talking, pressing a "push to talk" button on their radio which switches off the receiver and switches on the transmitter. Or the radio link may be full duplex, a bidirectional link using two radio channels so both people can talk at the same time, as in a cell phone.
Cell phone – a portable wireless telephone that is connected to the telephone network by radio signals exchanged with a local antenna at a cellular base station (cell tower). The service area covered by the provider is divided into small geographical areas called "cells", each served by a separate base station antenna and multichannel transceiver. All the cell phones in a cell communicate with this antenna on separate frequency channels, assigned from a common pool of frequencies. The purpose of cellular organization is to conserve radio bandwidth by frequency reuse. Low power transmitters are used so the radio waves used in a cell do not travel far beyond the cell, allowing the same frequencies to be reused in geographically separated cells. When a user carrying a cellphone crosses from one cell to another, his phone is automatically "handed off" seamlessly to the new antenna and assigned new frequencies. Cellphones have a highly automated full duplex digital transceiver using OFDM modulation using two digital radio channels, each carrying one direction of the bidirectional conversation, as well as a control channel that handles dialing calls and "handing off" the phone to another cell tower. Older 2G, 3G, and 4G networks use frequencies in the UHF and low microwave range, between 700 MHz and 3 GHz. The cell phone transmitter adjusts its power output to use the minimum power necessary to communicate with the cell tower; 0.6 W when near the tower, up to 3 W when farther away. Cell tower channel transmitter power is 50 W. Current generation phones, called smartphones, have many functions besides making telephone calls, and therefore have several other radio transmitters and receivers that connect them with other networks: usually a Wi-Fi modem, a Bluetooth modem, and a GPS receiver.
5G cellular network – next-generation cellular networks which began deployment in 2019. Their major advantage is much higher data rates than previous cellular networks, up to 10 Gbps; 100 times faster than the previous cellular technology, 4G LTE. The higher data rates are achieved partly by using higher frequency radio waves, in the higher microwave band 3–6 GHz, and millimeter wave band, around 28 and 39 GHz. Since these frequencies have a shorter range than previous cellphone bands, the cells will be smaller than the cells in previous cellular networks which could be many miles across. Millimeter-wave cells will only be a few blocks long, and instead of a cell base station and antenna tower, they will have many small antennas attached to utility poles and buildings.
Satellite phone (satphone) – a portable wireless telephone similar to a cell phone, connected to the telephone network through a radio link to an orbiting communications satellite instead of through cell towers. They are more expensive than cell phones; but their advantage is that, unlike a cell phone which is limited to areas covered by cell towers, satphones can be used over most or all of the geographical area of the Earth. In order for the phone to communicate with a satellite using a small omnidirectional antenna, first-generation systems use satellites in low Earth orbit, about above the surface. With an orbital period of about 100 minutes, a satellite can only be in view of a phone for about 4 – 15 minutes, so the call is "handed off" to another satellite when one passes beyond the local horizon. Therefore, large numbers of satellites, about 40 to 70, are required to ensure that at least one satellite is in view continuously from each point on Earth. Other satphone systems use satellites in geostationary orbit in which only a few satellites are needed, but these cannot be used at high latitudes because of terrestrial interference.
Cordless phone – a landline telephone in which the handset is portable and communicates with the rest of the phone by a short-range full duplex radio link, instead of being attached by a cord. Both the handset and the base station have low-power radio transceivers that handle the short-range bidirectional radio link. , cordless phones in most nations use the DECT transmission standard.
Land mobile radio system – short-range mobile or portable half-duplex radio transceivers operating in the VHF or UHF band that can be used without a license. They are often installed in vehicles, with the mobile units communicating with a dispatcher at a fixed base station. Special systems with reserved frequencies are used by first responder services; police, fire, ambulance, and emergency services, and other government services. Other systems are made for use by commercial firms such as taxi and delivery services. VHF systems use channels in the range 30–50 MHz and 150–172 MHz. UHF systems use the 450–470 MHz band and in some areas the 470–512 MHz range. In general, VHF systems have a longer range than UHF but require longer antennas. AM or FM modulation is mainly used, but digital systems such as DMR are being introduced. The radiated power is typically limited to 4 watts. These systems have a fairly limited range, usually depending on terrain. Repeaters installed on tall buildings, hills, or mountain peaks are often used to increase the range when it is desired to cover a larger area than line-of-sight. Examples of land mobile systems are CB, FRS, GMRS, and MURS. Modern digital systems, called trunked radio systems, have a digital channel management system using a control channel that automatically assigns frequency channels to user groups.
Walkie-talkie – a battery-powered portable handheld half-duplex two-way radio, used in land mobile radio systems.
Airband – Half-duplex radio system used by aircraft pilots to talk to other aircraft and ground-based air traffic controllers. This vital system is the main communication channel for air traffic control. For most communication in overland flights in air corridors a VHF-AM system using channels between 108 and 137 MHz in the VHF band is used. This system has a typical transmission range of for aircraft flying at cruising altitude. For flights in more remote areas, such as transoceanic airline flights, aircraft use the HF band or channels on the Inmarsat or Iridium satphone satellites. Military aircraft also use a dedicated UHF-AM band from 225.0 to 399.95 MHz.
Marine radio – medium-range transceivers on ships, used for ship-to-ship, ship-to-air, and ship-to-shore communication with harbormasters They use FM channels between 156 and 174 MHz in the VHF band with up to 25 watts power, giving them a range of about . Some channels are half-duplex and some are full-duplex, to be compatible with the telephone network, to allow users to make telephone calls through a marine operator.
Amateur radio – long-range half-duplex two-way radio used by hobbyists for non-commercial purposes: recreational radio contacts with other amateurs, volunteer emergency communication during disasters, contests, and experimentation. Radio amateurs must hold an amateur radio license and are given a unique callsign that must be used as an identifier in transmissions. Amateur radio is restricted to small frequency bands, the amateur radio bands, spaced throughout the radio spectrum starting at 136 kHz. Within these bands, amateurs are allowed the freedom to transmit on any frequency using a wide variety of voice modulation methods, along with other forms of communication, such as slow-scan television (SSTV), and radioteletype (RTTY). Additionally, amateurs are among the only radio operators still using Morse code radiotelegraphy.
One-way voice communication
One way, unidirectional radio transmission is called simplex.
Baby monitor – a crib-side appliance for parents of infants that transmits the baby's sounds to a receiver carried by the parent, so they can monitor the baby while they are in other parts of the house. The wavebands used vary by region, but analog baby monitors generally transmit with low power in the 16, 9.3–49.9 or 900 MHz wavebands, and digital systems in the 2.4 GHz waveband. Many baby monitors have duplex channels so the parent can talk to the baby, and cameras to show video of the baby.
Wireless microphone – a battery-powered microphone with a short-range transmitter that is handheld or worn on a person's body which transmits its sound by radio to a nearby receiver unit connected to a sound system. Wireless microphones are used by public speakers, performers, and television personalities so they can move freely without trailing a microphone cord. Traditionally, analog models transmit in FM on unused portions of the television broadcast frequencies in the VHF and UHF bands. Some models transmit on two frequency channels for diversity reception to prevent nulls from interrupting transmission as the performer moves around. Some models use digital modulation to prevent unauthorized reception by scanner radio receivers; these operate in the 900 MHz, 2.4 GHz or 6 GHz ISM bands. European standards also support wireless multichannel audio systems (WMAS) that can better support the use of large numbers of wireless microphones at a single event or venue. , U.S. regulators were considering adopting rules for WMAS.
Data communication
Wireless networking – automated radio links which transmit digital data between computers and other wireless devices using radio waves, linking the devices together transparently in a computer network. Computer networks can transmit any form of data: in addition to email and web pages, they also carry phone calls (VoIP), audio, and video content (called streaming media). Security is more of an issue for wireless networks than for wired networks since anyone nearby with a wireless modem can access the signal and attempt to log in. The radio signals of wireless networks are encrypted using WPA.
Wireless LAN (wireless local area network or Wi-Fi) – based on the IEEE 802.11 standards, these are the most widely used computer networks, used to implement local area networks without cables, linking computers, laptops, cell phones, video game consoles, smart TVs and printers in a home or office together, and to a wireless router connecting them to the Internet with a wire or cable connection. Wireless routers in public places like libraries, hotels and coffee shops create wireless access points (hotspots) to allow the public to access the Internet with portable devices like smartphones, tablets or laptops. Each device exchanges data using a wireless modem (wireless network interface controller), an automated microwave transmitter and receiver with an omnidirectional antenna that works in the background, exchanging data packets with the router. Wi-Fi uses channels in the 2.4 GHz and 5 GHz ISM bands with OFDM (orthogonal frequency-division multiplexing) modulation to transmit data at high rates. The transmitters in Wi-Fi modems are limited to a radiated power of 200 mW to 1 watt, depending on country. They have a maximum indoor range of about on 2.4 GHz and on 5 GHz.
Wireless WAN (wireless wide area network, WWAN) – a variety of technologies that provide wireless internet access over a wider area than Wi-Fi networks do – from an office building to a campus to a neighborhood, or to an entire city. The most common technologies used are: cellular modems, that exchange computer data by radio with cell towers; satellite internet access; and lower frequencies in the UHF band, which have a longer range than Wi-Fi frequencies. Since WWAN networks are much more expensive and complicated to administer than Wi-Fi networks, their use so far has generally been limited to private networks operated by large corporations.
Bluetooth – a very short-range wireless interface on a portable wireless device used as a substitute for a wire or cable connection, mainly to exchange files between portable devices and connect cellphones and music players with wireless headphones. In the most widely used mode, transmission power is limited to 1 milliwatt, giving it a very short range of up to 10 m (30 feet). The system uses frequency-hopping spread spectrum transmission, in which successive data packets are transmitted in a pseudorandom order on one of 79 1 MHz Bluetooth channels between 2.4 and 2.83 GHz in the ISM band. This allows Bluetooth networks to operate in the presence of noise, other wireless devices and other Bluetooth networks using the same frequencies, since the chance of another device attempting to transmit on the same frequency at the same time as the Bluetooth modem is low. In the case of such a "collision", the Bluetooth modem just retransmits the data packet on another frequency.
Packet radio – a long-distance peer-to-peer wireless ad-hoc network in which data packets are exchanged between computer-controlled radio modems (transmitter/receivers) called nodes, which may be separated by miles, and maybe mobile. Each node only communicates with neighboring nodes, so packets of data are passed from node to node until they reach their destination using the X.25 network protocol. Packet radio systems are used to a limited degree by commercial telecommunications companies and by the amateur radio community.
Text messaging (texting) – this is a service on cell phones, allowing a user to type a short alphanumeric message and send it to another phone number, and the text is displayed on the recipient's phone screen. It is based on the Short Message Service (SMS) which transmits using spare bandwidth on the control radio channel used by cell phones to handle background functions like dialing and cell handoffs. Due to technical limitations of the channel, text messages are limited to 160 alphanumeric characters.
Microwave relay – a long-distance high bandwidth point-to-point digital data transmission link consisting of a microwave transmitter connected to a dish antenna that transmits a beam of microwaves to another dish antenna and receiver. Since the antennas must be in line-of-sight, distances are limited by the visual horizon to . Microwave links are used for private business data, wide area computer networks (WANs), and by telephone companies to transmit long-distance phone calls and television signals between cities.
Telemetry – automated one-way (simplex) transmission of measurements and operation data from a remote process or device to a receiver for monitoring. Telemetry is used for in-flight monitoring of missiles, drones, satellites, and weather balloon radiosondes, sending scientific data back to Earth from interplanetary spacecraft, communicating with electronic biomedical sensors implanted in the human body, and well logging. Multiple channels of data are often transmitted using frequency-division multiplexing or time-division multiplexing. Telemetry is starting to be used in consumer applications such as:
Automated meter reading – electric power meters, water meters, and gas meters that, when triggered by an interrogation signal, transmit their readings by radio to a utility reader vehicle at the curb, to eliminate the need for an employee to go on the customer's property to manually read the meter.
Electronic toll collection – on toll roads, an alternative to manual collection of tolls at a toll booth, in which a transponder in a vehicle, when triggered by a roadside transmitter, transmits a signal to a roadside receiver to register the vehicle's use of the road, enabling the owner to be billed for the toll.
Radio Frequency Identification (RFID) – identification tags containing a tiny radio transponder (receiver and transmitter) which are attached to merchandise. When it receives an interrogation pulse of radio waves from a nearby reader unit, the tag transmits back an ID number, which can be used to inventory goods. Passive tags, the most common type, have a chip powered by the radio energy received from the reader, rectified by a diode, and can be as small as a grain of rice. They are incorporated in products, clothes, railroad cars, library books, airline baggage tags and are implanted under the skin in pets and livestock (microchip implant) and even people. Privacy concerns have been addressed with tags that use encrypted signals and authenticate the reader before responding. Passive tags use 125–134 kHz, 13, 900 MHz and 2.4 and 5 GHz ISM bands and have a short range. Active tags, powered by a battery, are larger but can transmit a stronger signal, giving them a range of hundreds of meters.
Submarine communication – When submerged, submarines are cut off from all ordinary radio communication with their military command authorities by the conductive seawater. However radio waves of low enough frequencies, in the VLF (30 to 3 kHz) and ELF (below 3 kHz) bands are able to penetrate seawater. Navies operate large shore transmitting stations with power output in the megawatt range to transmit encrypted messages to their submarines in the world's oceans. Due to the small bandwidth, these systems cannot transmit voice, only text messages at a slow data rate. The communication channel is one-way, since the long antennas needed to transmit VLF or ELF waves cannot fit on a submarine. VLF transmitters use miles long wire antennas like umbrella antennas. A few nations use ELF transmitters operating around 80 Hz, which can communicate with submarines at lower depths. These use even larger antennas called ground dipoles, consisting of two ground (Earth) connections apart, linked by overhead transmission lines to a power plant transmitter.
Space communication
This is radio communication between a spacecraft and an Earth-based ground station, or another spacecraft. Communication with spacecraft involves the longest transmission distances of any radio links, up to billions of kilometers for interplanetary spacecraft. In order to receive the weak signals from distant spacecraft, satellite ground stations use large parabolic "dish" antennas up to in diameter and extremely sensitive receivers. High frequencies in the microwave band are used, since microwaves pass through the ionosphere without refraction, and at microwave frequencies the high-gain antennas needed to focus the radio energy into a narrow beam pointed at the receiver are small and take up a minimum of space in a satellite. Portions of the UHF, L, C, S, ku and ka band are allocated for space communication. A radio link that transmits data from the Earth's surface to a spacecraft is called an uplink, while a link that transmits data from the spacecraft to the ground is called a downlink.
Communication satellite – an artificial satellite used as a telecommunications relay to transmit data between widely separated points on Earth. These are used because the microwaves used for telecommunications travel by line of sight and so cannot propagate around the curve of the Earth. , there were 2,224 communications satellites in Earth orbit. Most are in geostationary orbit above the equator, so that the satellite appears stationary at the same point in the sky, so the satellite dish antennas of ground stations can be aimed permanently at that spot and do not have to move to track it. In a satellite ground station a microwave transmitter and large satellite dish antenna transmit a microwave uplink beam to the satellite. The uplink signal carries many channels of telecommunications traffic, such as long-distance telephone calls, television programs, and internet signals, using a technique called frequency-division multiplexing (FDM). On the satellite, a transponder receives the signal, translates it to a different downlink frequency to avoid interfering with the uplink signal, and retransmits it down to another ground station, which may be widely separated from the first. There the downlink signal is demodulated and the telecommunications traffic it carries is sent to its local destinations through landlines. Communication satellites typically have several dozen transponders on different frequencies, which are leased by different users.
Direct broadcast satellite – a geostationary communication satellite that transmits retail programming directly to receivers in subscriber's homes and vehicles on Earth, in satellite radio and TV systems. It uses a higher transmitter power than other communication satellites, to allow the signal to be received by consumers with a small unobtrusive antenna. For example, satellite television uses downlink frequencies from 12.2 to 12.7 GHz in the ku band transmitted at 100 to 250 watts, which can be received by relatively small satellite dishes mounted on the outside of buildings.
Other applications
Radar
Radar is a radiolocation method used to locate and track aircraft, spacecraft, missiles, ships, vehicles, and also to map weather patterns and terrain. A radar set consists of a transmitter and receiver. The transmitter emits a narrow beam of radio waves which is swept around the surrounding space. When the beam strikes a target object, radio waves are reflected back to the receiver. The direction of the beam reveals the object's location. Since radio waves travel at a constant speed close to the speed of light, by measuring the brief time delay between the outgoing pulse and the received "echo", the range to the target can be calculated. The targets are often displayed graphically on a map display called a radar screen. Doppler radar can measure a moving object's velocity, by measuring the change in frequency of the return radio waves due to the Doppler effect.
Radar sets mainly use high frequencies in the microwave bands, because these frequencies create strong reflections from objects the size of vehicles and can be focused into narrow beams with compact antennas. Parabolic (dish) antennas are widely used. In most radars the transmitting antenna also serves as the receiving antenna; this is called a monostatic radar. A radar which uses separate transmitting and receiving antennas is called a bistatic radar.
Airport surveillance radar – In aviation, radar is the main tool of air traffic control. A rotating dish antenna sweeps a vertical fan-shaped beam of microwaves around the airspace and the radar set shows the location of aircraft as "blips" of light on a display called a radar screen. Airport radar operates at 2.7 – 2.9 GHz in the microwave S band. In large airports the radar image is displayed on multiple screens in an operations room called the TRACON (Terminal Radar Approach Control), where air traffic controllers direct the aircraft by radio to maintain safe aircraft separation.
Secondary surveillance radar – Aircraft carry radar transponders, transceivers which when triggered by the incoming radar signal transmit a return microwave signal. This causes the aircraft to show up more strongly on the radar screen. The radar which triggers the transponder and receives the return beam, usually mounted on top of the primary radar dish, is called the secondary surveillance radar. Since radar cannot measure an aircraft's altitude with any accuracy, the transponder also transmits back the aircraft's altitude measured by its altimeter, and an ID number identifying the aircraft, which is displayed on the radar screen.
Electronic countermeasures (ECM) – Military defensive electronic systems designed to degrade enemy radar effectiveness, or deceive it with false information, to prevent enemies from locating local forces. It often consists of powerful microwave transmitters that can mimic enemy radar signals to create false target indications on the enemy radar screens.
Marine radar – an S or X band radar on ships used to detect nearby ships and obstructions like bridges. A rotating antenna sweeps a vertical fan-shaped beam of microwaves around the water surface surrounding the craft out to the horizon.
Weather radar – A Doppler radar which maps weather precipitation intensities and wind speeds with the echoes returned from raindrops and their radial velocity by their Doppler shift.
Phased-array radar – a radar set that uses a phased array, a computer-controlled antenna that can steer the radar beam quickly to point in different directions without moving the antenna. Phased-array radars were developed by the military to track fast-moving missiles and aircraft. They are widely used in military equipment and are now spreading to civilian applications.
Synthetic aperture radar (SAR) – a specialized airborne radar set that produces a high-resolution map of ground terrain. The radar is mounted on an aircraft or spacecraft and the radar antenna radiates a beam of radio waves sideways at right angles to the direction of motion, toward the ground. In processing the return radar signal, the motion of the vehicle is used to simulate a large antenna, giving the radar a higher resolution.
Ground-penetrating radar – a specialized radar instrument that is rolled along the ground surface in a cart and transmits a beam of radio waves into the ground, producing an image of subsurface objects. Frequencies from 100 MHz to a few GHz are used. Since radio waves cannot penetrate very far into earth, the depth of GPR is limited to about 50 feet.
Collision avoidance system – a short range radar or LIDAR system on an automobile or vehicle that detects if the vehicle is about to collide with an object and applies the brakes to prevent the collision.
Radar fuze – a detonator for an aerial bomb which uses a radar altimeter to measure the height of the bomb above the ground as it falls and detonates it at a certain altitude.
Radiolocation
Radiolocation is a generic term covering a variety of techniques that use radio waves to find the location of objects, or for navigation.
Global Navigation Satellite System (GNSS) or satnav system – A system of satellites which allows geographical location on Earth (latitude, longitude, and altitude/elevation) to be determined to high precision (within a few metres) by small portable navigation instruments, by timing the arrival of radio signals from the satellites. These are the most widely used navigation systems today. The main satellite navigation systems are the US Global Positioning System (GPS), Russia's GLONASS, China's BeiDou Navigation Satellite System (BDS) and the European Union's Galileo.
Global Positioning System (GPS) – The most widely used satellite navigation system, maintained by the US Air Force, which uses a constellation of 31 satellites in low Earth orbit. The orbits of the satellites are distributed so at any time at least four satellites are above the horizon over each point on Earth. Each satellite has an onboard atomic clock and transmits a continuous radio signal containing a precise time signal as well as its current position. Two frequencies are used, 1.2276 and 1.57542 GHz. Since the velocity of radio waves is virtually constant, the delay of the radio signal from a satellite is proportional to the distance of the receiver from the satellite. By receiving the signals from at least four satellites a GPS receiver can calculate its position on Earth by comparing the arrival time of the radio signals. Since each satellite's position is known precisely at any given time, from the delay the position of the receiver can be calculated by a microprocessor in the receiver. The position can be displayed as latitude and longitude, or as a marker on an electronic map. GPS receivers are incorporated in almost all cellphones and in vehicles such as automobiles, aircraft, and ships, and are used to guide drones, missiles, cruise missiles, and even artillery shells to their target, and handheld GPS receivers are produced for hikers and the military.
Radio beacon – a fixed location terrestrial radio transmitter which transmits a continuous radio signal used by aircraft and ships for navigation. The locations of beacons are plotted on navigational maps used by aircraft and ships.
VHF omnidirectional range (VOR) – a worldwide aircraft radio navigation system consisting of fixed ground radio beacons transmitting between 108.00 and 117.95 MHz in the very high frequency (VHF) band. An automated navigational instrument on the aircraft displays a bearing to a nearby VOR transmitter. A VOR beacon transmits two signals simultaneously on different frequencies. A directional antenna transmits a beam of radio waves that rotates like a lighthouse at a fixed rate, 30 times per second. When the directional beam is facing north, an omnidirectional antenna transmits a pulse. By measuring the difference in phase of these two signals, an aircraft can determine its bearing (or "radial") from the station accurately. By taking a bearing on two VOR beacons an aircraft can determine its position (called a "fix") to an accuracy of about . Most VOR beacons also have a distance measuring capability, called distance measuring equipment (DME); these are called VOR/DME's. The aircraft transmits a radio signal to the VOR/DME beacon and a transponder transmits a return signal. From the propagation delay between the transmitted and received signal the aircraft can calculate its distance from the beacon. This allows an aircraft to determine its location "fix" from only one VOR beacon. Since line-of-sight VHF frequencies are used VOR beacons have a range of about 200 miles for aircraft at cruising altitude. TACAN is a similar military radio beacon system which transmits in 962–1213 MHz, and a combined VOR and TACAN beacon is called a VORTAC. The number of VOR beacons is declining as aviation switches to the RNAV system that relies on Global Positioning System satellite navigation.
Instrument Landing System (ILS) - A short range radio navigation aid at airports which guides aircraft landing in low visibility conditions. It consists of multiple antennas at the end of each runway that radiate two beams of radio waves along the approach to the runway: the localizer (108 to 111.95 MHz frequency), which provides horizontal guidance, a heading line to keep the aircraft centered on the runway, and the glideslope (329.15 to 335 MHz) for vertical guidance, to keep the aircraft descending at the proper rate for a smooth touchdown at the correct point on the runway. Each aircraft has a receiver instrument and antenna which receives the beams, with an indicator to tell the pilot whether he is on the correct horizontal and vertical approach. The ILS beams are receivable for at least 15 miles, and have a radiated power of 25 watts. ILS systems at airports are being replaced by systems that use satellite navigation.
Non-directional beacon (NDB) – Legacy fixed radio beacons used before the VOR system that transmit a simple signal in all directions for aircraft or ships to use for radio direction finding. Aircraft use automatic direction finder (ADF) receivers which use a directional antenna to determine the bearing to the beacon. By taking bearings on two beacons they can determine their position. NDBs use frequencies between 190 and 1750 kHz in the LF and MF bands which propagate beyond the horizon as ground waves or skywaves much farther than VOR beacons. They transmit a callsign consisting of one to 3 Morse code letters as an identifier.
Emergency locator beacon – a portable battery powered radio transmitter used in emergencies to locate airplanes, vessels, and persons in distress and in need of immediate rescue. Various types of emergency locator beacons are carried by aircraft, ships, vehicles, hikers and cross-country skiers. In the event of an emergency, such as the aircraft crashing, the ship sinking, or a hiker becoming lost, the transmitter is deployed and begins to transmit a continuous radio signal, which is used by search and rescue teams to quickly find the emergency and render aid. The latest generation Emergency Position Indicating Rescue Beacons (EPIRBs) contain a GPS receiver, and broadcast to rescue teams their exact location within 20 meters.
Cospas-Sarsat – an international humanitarian consortium of governmental and private agencies which acts as a dispatcher for search and rescue operations. It operates a network of about 47 satellites carrying radio receivers, which detect distress signals from emergency locator beacons anywhere on Earth transmitting on the international Cospas distress frequency of 406 MHz. The satellites calculate the geographic location of the beacon within 2 km by measuring the Doppler frequency shift of the radio waves due to the relative motion of the transmitter and the satellite, and quickly transmit the information to the appropriate local first responder organizations, which perform the search and rescue.
Radio direction finding (RDF) – this is a general technique, used since the early 1900s, of using specialized radio receivers with directional antennas (RDF receivers) to determine the exact bearing of a radio signal, to determine the location of the transmitter. The location of a terrestrial transmitter can be determined by simple triangulation from bearings taken by two RDF stations separated geographically, as the point where the two bearing lines cross, this is called a "fix". Military forces use RDF to locate enemy forces by their tactical radio transmissions, counterintelligence services use it to locate clandestine transmitters used by espionage agents, and governments use it to locate unlicensed transmitters or interference sources. Older RDF receivers used rotatable loop antennas, the antenna is rotated until the radio signal strength is weakest, indicating the transmitter is in one of the antenna's two nulls. The nulls are used since they are sharper than the antenna's lobes (maxima). More modern receivers use phased array antennas which have a much greater angular resolution.
Animal migration tracking – a widely used technique in wildlife biology, conservation biology, and wildlife management in which small battery-powered radio transmitters are attached to wild animals so their movements can be tracked with a directional RDF receiver. Sometimes the transmitter is implanted in the animal. The VHF band is typically used since antennas in this band are fairly compact. The receiver has a directional antenna (typically a small Yagi) which is rotated until the received signal is strongest; at this point the antenna is pointing in the direction of the animal. Sophisticated systems used in recent years use satellites to track the animal, or geolocation tags with GPS receivers which record and transmit a log of the animal's location.
Remote control
Radio remote control is the use of electronic control signals sent by radio waves from a transmitter to control the actions of a device at a remote location. Remote control systems may also include telemetry channels in the other direction, used to transmit real-time information on the state of the device back to the control station. Uncrewed spacecraft are an example of remote-controlled machines, controlled by commands transmitted by satellite ground stations. Most handheld remote controls used to control consumer electronics products like televisions or DVD players actually operate by infrared light rather than radio waves, so are not examples of radio remote control. A security concern with remote control systems is spoofing, in which an unauthorized person transmits an imitation of the control signal to take control of the device. Examples of radio remote control:
Unmanned aerial vehicle (UAV, drone) – A drone is an aircraft without an onboard pilot, flown by remote control by a pilot in another location, usually in a piloting station on the ground. They are used by the military for reconnaissance and ground attack, and more recently by the civilian world for news reporting and aerial photography. The pilot uses aircraft controls like a joystick or steering wheel, which create control signals which are transmitted to the drone by radio to control the flight surfaces and engine. A telemetry system transmits back a video image from a camera in the drone to allow the pilot to see where the aircraft is going, and data from a GPS receiver giving the real-time position of the aircraft. UAVs have sophisticated onboard automatic pilot systems that maintain stable flight and only require manual control to change directions.
Keyless entry system – a short-range handheld battery powered key fob transmitter, included with most modern cars, which can lock and unlock the doors of a vehicle from outside, eliminating the need to use a key. When a button is pressed, the transmitter sends a coded radio signal to a receiver in the vehicle, operating the locks. The fob must be close to the vehicle, typically within 5 to 20 meters. North America and Japan use a frequency of 315 MHz, while Europe uses 433.92 and 868 MHz. Some models can also remotely start the engine, to warm up the car. A security concern with all keyless entry systems is a replay attack, in which a thief uses a special receiver ("code grabber") to record the radio signal during opening, which can later be replayed to open the door. To prevent this, keyless systems use a rolling code system in which a pseudorandom number generator in the remote control generates a different random key each time it is used. To prevent thieves from simulating the pseudorandom generator to calculate the next key, the radio signal is also encrypted.
Garage door opener – a short-range handheld transmitter which can open or close a building's electrically operated garage door from outside, so the owner can open the door upon arrival, and close it after departure. When a button is pressed the control transmits a coded FSK radio signal to a receiver in the opener, raising or lowering the door. Modern openers use 310, 315 or 390 MHz. To prevent a thief using a replay attack, modern openers use a rolling code system.
Radio-controlled models – a popular hobby is playing with radio-controlled model boats, cars, airplanes, and helicopters (quadcopters) which are controlled by radio signals from a handheld console with a joystick. Most recent transmitters use the 2.4 GHz ISM band with multiple control channels modulated with PWM, PCM or FSK.
Wireless doorbell – A residential doorbell that uses wireless technology to eliminate the need to run wires through the building walls. It consists of a doorbell button beside the door containing a small battery powered transmitter. When the doorbell is pressed it sends a signal to a receiver inside the house with a speaker that sounds chimes to indicate someone is at the door. They usually use the 2.4 GHz ISM band. The frequency channel used can usually be changed by the owner in case another nearby doorbell is using the same channel.
Scientific research
Radio astronomy is the scientific study of radio waves emitted by astronomical objects. Radio astronomers use radio telescopes, large radio antennas and receivers, to receive and study the radio waves from astronomical radio sources. Since astronomical radio sources are so far away, the radio waves from them are extremely weak, requiring extremely sensitive receivers, and radio telescopes are the most sensitive radio receivers in existence. They use large parabolic (dish) antennas up to in diameter to collect enough radio wave energy to study. The RF front end electronics of the receiver is often cooled by liquid nitrogen to reduce thermal noise. Multiple antennas are often linked together in arrays which function as a single antenna, to increase collecting power. In Very Long Baseline Interferometry (VLBI) radio telescopes on different continents are linked, which can achieve the resolution of an antenna thousands of miles in diameter.
Remote sensing – in radio, remote sensing is the reception of electromagnetic waves radiated by natural objects or the atmosphere for scientific research. All warm objects emit microwaves and the spectrum emitted can be used to determine temperature. Microwave radiometers are used in meteorology and earth sciences to determine temperature of the atmosphere and earth surface, as well as chemical reactions in the atmosphere.
Jamming
Radio jamming is the deliberate radiation of radio signals designed to interfere with the reception of other radio signals. Jamming devices are called "signal suppressors" or "interference generators" or just jammers.
During wartime, militaries use jamming to interfere with enemies' tactical radio communication. Since radio waves can pass beyond national borders, some totalitarian countries which practice censorship use jamming to prevent their citizens from listening to broadcasts from radio stations in other countries. Jamming is usually accomplished by a powerful transmitter which generates noise on the same frequency as the target transmitter.
US Federal law prohibits the nonmilitary operation or sale of any type of jamming devices, including ones that interfere with GPS, cellular, Wi-Fi and police radars.
| Technology | Media and communication | null |
19304813 | https://en.wikipedia.org/wiki/Burnetia | Burnetia | Burnetia is an extinct genus of biarmosuchian therapsids in the family Burnetiidae, from the Late Permian of South Africa. Burnetia is known so far from a single holotype skull lacking the lower jaws described by South African paleontologist Robert Broom in 1923. Due to erosion and dorsoventral crushing, features of the skull are hard to interpret. Stutural lines are further distorted by the unusual shape of the skull roof, including many bosses and protuberances.
Description
When broadly looking at the skull, there are well-developed "cheeks", bosses and pits that resemble Pareiasaurians'. However, the small temporal fossa distinguishes it from the Cotylosaur. The overall shape resembles a triangle. In the nasals, there is a bulging expansion of bone. Unlike proburnetia's median nasal bridge being long, narrow and raised, Burnetia is splindle-shaped. The median nasal boss is spindle-shaped. The snout is wide and blunt. The large preorbital pits on the lachrymal are significant. Over the orbit there are notable ridges on the prefrontal and frontal. The supra-orbital ridges make the orbits face distally and posteriorly. The suborbital eminence is subdivided into distinguished portions. The small pineal foramen sits dorsally on a boss.
Burnetia palate is similar to Gorgonopsians'. Anteriorly, the internal nares have the lower canines. The maxilla is adjacent to the palatine, and posterior to the palatine is the long prevomer that meets the premaxilla. The palatine is tooth bearing, as well as the pterygoid that is just posterior to the palatine. The vomer is held by the surrounding vomerine processes that form the choanae's middle border. Unlike the rest of burnetiamorphs, Burnetia interchoanal part of the vomer is not narrow.
The concave occiput is tilted up, which is shown when it is aligned vertically, the snout faces downward. The supraoccipital sits anterior to the paraoccipital. The size of the basioccipital is considered to be small. The majority of the occipital surface, posterior "horn", and posterior lateral margins are made from the squamosal.
The basisphenoid in Burnetia differs from Gorgonopsians'. In Burnetia, their basisphenoid is round and shallow. In the middle of the basisphenoid, it is separated by a groove. Gorgonopsians' basisphenoid contrasts Burnetia by having a "single sharp median keel".
Discovery
Broom found Burnetia mirabilis in the Dicynodon Assemblage Zone near Graaff-Reinet, South Africa. Broom concluded that Burnetia was related closest to the group of Gorgonopsians. However, Broom made observations based on the skull when it was covered by matrix and no underlying bone was visible. Lieuwe Boonstra removed this matrix and found that Burnetia and Gorgonopsians differed mainly by the thickening of dermal bones, bosses and their basisphenoids.
Paleoenvironment
Phylogenetic analysis done by Ashley Kruger suggests that a likely origin for burnetiamorphs could be southern Africa. The base and oldest burnetiamorphs are found in South Africa.
Classification and related taxa
The family Burnetiidae came from the discovery of Burnetia, but new discoveries led to the studies of burnetiamorphs. Burnetiamorphs only have about two taxa that represent each genus. It was previously believed that the Burnetiidae clade only contained two taxa, Burnetia mirabilis and Proburnetia viatkensis, but later Pachydectes and Bullacephalus were found to be included in the clade, as well. It was unable to be confirmed if Pachydectes shared the feature of nasal the nasal supporting a middle boss like Burnetia. However, like Burnetia, above the postorbital bar they have a significant boss. Lycaenodon longiceps is in the clade Biarmosuchia and has some similarities to Burnetia. Both Lycaenodon and Burnetia likely had a large foramen, surrounded by the vomer and vomerine process. This foramen implies that in biarmosuchians, the vomeronasal organ may have communicated with the oral cavity. They also share the trait of having a long vomerine process. The first burnetiid therapsid found in Tanzania's Usili Formation was found to most resemble Burnetia. The Usili burnetiid and Burnetia both had bosses above their orbits that were "S"-like.
| Biology and health sciences | Proto-mammals | Animals |
6026198 | https://en.wikipedia.org/wiki/Monty%20Hall%20problem | Monty Hall problem | The Monty Hall problem is a brain teaser, in the form of a probability puzzle, based nominally on the American television game show Let's Make a Deal and named after its original host, Monty Hall. The problem was originally posed (and solved) in a letter by Steve Selvin to the American Statistician in 1975. It became famous as a question from reader Craig F. Whitaker's letter quoted in Marilyn vos Savant's "Ask Marilyn" column in Parade magazine in 1990:
Savant's response was that the contestant should switch to the other door. By the standard assumptions, the switching strategy has a probability of winning the car, while the strategy of keeping the initial choice has only a probability.
When the player first makes their choice, there is a chance that the car is behind one of the doors not chosen. This probability does not change after the host reveals a goat behind one of the unchosen doors. When the host provides information about the two unchosen doors (revealing that one of them does not have the car behind it), the chance of the car being behind one of the unchosen doors rests on the unchosen and unrevealed door, as opposed to the chance of the car being behind the door the contestant chose initially.
The given probabilities depend on specific assumptions about how the host and contestant choose their doors. An important insight is that, with these standard conditions, there is more information about doors 2 and 3 than was available at the beginning of the game when door 1 was chosen by the player: the host's action adds value to the door not eliminated, but not to the one chosen by the contestant originally. Another insight is that switching doors is a different action from choosing between the two remaining doors at random, as the former action uses the previous information and the latter does not. Other possible behaviors of the host than the one described can reveal different additional information, or none at all, leading to different probabilities. In her response, Savant states:
Many readers of Savant's column refused to believe switching is beneficial and rejected her explanation. After the problem appeared in Parade, approximately 10,000 readers, including nearly 1,000 with PhDs, wrote to the magazine, most of them calling Savant wrong. Even when given explanations, simulations, and formal mathematical proofs, many people still did not accept that switching is the best strategy. Paul Erdős, one of the most prolific mathematicians in history, remained unconvinced until he was shown a computer simulation demonstrating Savant's predicted result.
The problem is a paradox of the veridical type, because the solution is so counterintuitive it can seem absurd but is nevertheless demonstrably true. The Monty Hall problem is mathematically related closely to the earlier three prisoners problem and to the much older Bertrand's box paradox.
Paradox
Steve Selvin wrote a letter to the American Statistician in 1975, describing a problem based on the game show Let's Make a Deal, dubbing it the "Monty Hall problem" in a subsequent letter. The problem is equivalent mathematically to the Three Prisoners problem described in Martin Gardner's "Mathematical Games" column in Scientific American in 1959 and the Three Shells Problem described in Gardner's book Aha Gotcha.
Standard assumptions
By the standard assumptions, the probability of winning the car after switching is .
This solution is due to the behavior of the host. Ambiguities in the Parade version do not explicitly define the protocol of the host. However, Marilyn vos Savant's solution printed alongside Whitaker's question implies, and both Selvin and Savant explicitly define, the role of the host as follows:
The host must always open a door that was not selected by the contestant.
The host must always open a door to reveal a goat and never the car.
The host must always offer the chance to switch between the door chosen originally and the closed door remaining.
When any of these assumptions is varied, it can change the probability of winning by switching doors as detailed in the section below. It is also typically presumed that the car is initially hidden randomly behind the doors and that, if the player initially chooses the car, then the host's choice of which goat-hiding door to open is random. Some authors, independently or inclusively, assume that the player's initial choice is random as well.
Simple solutions
The solution presented by Savant in Parade shows the three possible arrangements of one car and two goats behind three doors and the result of staying or switching after initially picking door 1 in each case:
A player who stays with the initial choice wins in only one out of three of these equally likely possibilities, while a player who switches wins in two out of three.
An intuitive explanation is that, if the contestant initially picks a goat (2 of 3 doors), the contestant will win the car by switching because the other goat can no longer be picked – the host had to reveal its location – whereas if the contestant initially picks the car (1 of 3 doors), the contestant will not win the car by switching. Using the switching strategy, winning or losing thus only depends on whether the contestant has initially chosen a goat ( probability) or the car ( probability). The fact that the host subsequently reveals a goat in one of the unchosen doors changes nothing about the initial probability.
Most people conclude that switching does not matter, because there would be a 50% chance of finding the car behind either of the two unopened doors. This would be true if the host selected a door to open at random, but this is not the case. The host-opened door depends on the player's initial choice, so the assumption of independence does not hold. Before the host opens a door, there is a probability that the car is behind each door. If the car is behind door 1, the host can open either door 2 or door 3, so the probability that the car is behind door 1 and the host opens door 3 is × = . If the car is behind door 2 – with the player having picked door 1 – the host must open door 3, such the probability that the car is behind door 2 and the host opens door 3 is × 1 = . These are the only cases where the host opens door 3, so if the player has picked door 1 and the host opens door 3, the car is twice as likely to be behind door 2 as door 1. The key is that if the car is behind door 2 the host must open door 3, but if the car is behind door 1 the host can open either door.
Another way to understand the solution is to consider together the two doors initially unchosen by the player. As Cecil Adams puts it, "Monty is saying in effect: you can keep your one door or you can have the other two doors". The chance of finding the car has not been changed by the opening of one of these doors because Monty, knowing the location of the car, is certain to reveal a goat. The player's choice after the host opens a door is no different than if the host offered the player the option to switch from the original chosen door to the set of both remaining doors. The switch in this case clearly gives the player a probability of choosing the car.
As Keith Devlin says, "By opening his door, Monty is saying to the contestant 'There are two doors you did not choose, and the probability that the prize is behind one of them is . I'll help you by using my knowledge of where the prize is to open one of those two doors to show you that it does not hide the prize. You can now take advantage of this additional information. Your choice of door A has a chance of 1 in 3 of being the winner. I have not changed that. But by eliminating door C, I have shown you that the probability that door B hides the prize is 2 in 3.
Savant suggests that the solution will be more intuitive with 1,000,000 doors rather than 3. In this case, there are 999,999 doors with goats behind them and one door with a prize. After the player picks a door, the host opens 999,998 of the remaining doors. On average, in 999,999 times out of 1,000,000, the remaining door will contain the prize. Intuitively, the player should ask how likely it is that, given a million doors, they managed to pick the right one initially. Stibel et al. proposed that working memory demand is taxed during the Monty Hall problem and that this forces people to "collapse" their choices into two equally probable options. They report that when the number of options is increased to more than 7 people tend to switch more often; however, most contestants still incorrectly judge the probability of success to be 50%.
Savant and the media furor
Savant wrote in her first column on the Monty Hall problem that the player should switch. She received thousands of letters from her readersthe vast majority of which, including many from readers with PhDs, disagreed with her answer. During 1990–1991, three more of her columns in Parade were devoted to the paradox. Numerous examples of letters from readers of Savant's columns are presented and discussed in The Monty Hall Dilemma: A Cognitive Illusion Par Excellence.
The discussion was replayed in other venues (e.g., in Cecil Adams' The Straight Dope newspaper column) and reported in major newspapers such as The New York Times.
In an attempt to clarify her answer, she proposed a shell game to illustrate: "You look away, and I put a pea under one of three shells. Then I ask you to put your finger on a shell. The odds that your choice contains a pea are , agreed? Then I simply lift up an empty shell from the remaining other two. As I can (and will) do this regardless of what you've chosen, we've learned nothing to allow us to revise the odds on the shell under your finger." She also proposed a similar simulation with three playing cards.
Savant commented that, though some confusion was caused by some readers' not realizing they were supposed to assume that the host must always reveal a goat, almost all her numerous correspondents had correctly understood the problem assumptions, and were still initially convinced that Savant's answer ("switch") was wrong.
Confusion and criticism
Sources of confusion
When first presented with the Monty Hall problem, an overwhelming majority of people assume that each door has an equal probability and conclude that switching does not matter. Out of 228 subjects in one study, only 13% chose to switch. In his book The Power of Logical Thinking, cognitive psychologist writes: "No other statistical puzzle comes so close to fooling all the people all the time [and] even Nobel physicists systematically give the wrong answer, and that they insist on it, and they are ready to berate in print those who propose the right answer". Pigeons repeatedly exposed to the problem show that they rapidly learn to always switch, unlike humans.
Most statements of the problem, notably the one in Parade, do not match the rules of the actual game show and do not fully specify the host's behavior or that the car's location is randomly selected. However, Krauss and Wang argue that people make the standard assumptions even if they are not explicitly stated.
Although these issues are mathematically significant, even when controlling for these factors, nearly all people still think each of the two unopened doors has an equal probability and conclude that switching does not matter. This "equal probability" assumption is a deeply rooted intuition. People strongly tend to think probability is evenly distributed across as many unknowns as are present, whether or not that is true in the particular situation under consideration.
The problem continues to attract the attention of cognitive psychologists. The typical behavior of the majority, i.e., not switching, may be explained by phenomena known in the psychological literature as:
The endowment effect, in which people tend to overvalue the winning probability of the door already chosenalready "owned".
The status quo bias, in which people prefer to keep the choice of door they have made already.
The errors of omission vs. errors of commission effect, in which, all other things being equal, people prefer to make errors by inaction (Stay) as opposed to action (Switch).
Experimental evidence confirms that these are plausible explanations that do not depend on probability intuition. Another possibility is that people's intuition simply does not deal with the textbook version of the problem, but with a real game show setting. There, the possibility exists that the show master plays deceitfully by opening other doors only if a door with the car was initially chosen. A show master playing deceitfully half of the times modifies the winning chances in case one is offered to switch to "equal probability".
Criticism of the simple solutions
As already remarked, most sources in the topic of probability, including many introductory probability textbooks, solve the problem by showing the conditional probabilities that the car is behind door 1 and door 2 are and (not and ) given that the contestant initially picks door 1 and the host opens door 3; various ways to derive and understand this result were given in the previous subsections.
Among these sources are several that explicitly criticize the popularly presented "simple" solutions, saying these solutions are "correct but ... shaky", or do not "address the problem posed", or are "incomplete", or are "unconvincing and misleading", or are (most bluntly) "false".
Sasha Volokh (2015) wrote that "any explanation that says something like 'the probability of door 1 was , and nothing can change that...' is automatically fishy: probabilities are expressions of our ignorance about the world, and new information can change the extent of our ignorance."
Some say that these solutions answer a slightly different questionone phrasing is "you have to announce before a door has been opened whether you plan to switch".
The simple solutions show in various ways that a contestant who is determined to switch will win the car with probability , and hence that switching is the winning strategy, if the player has to choose in advance between "always switching", and "always staying". However, the probability of winning by always switching is a logically distinct concept from the probability of winning by switching given that the player has picked door 1 and the host has opened door 3. As one source says, "the distinction between [these questions] seems to confound many". The fact that these are different can be shown by varying the problem so that these two probabilities have different numeric values. For example, assume the contestant knows that Monty does not open the second door randomly among all legal alternatives but instead, when given an opportunity to choose between two losing doors, Monty will open the one on the right. In this situation, the following two questions have different answers:
What is the probability of winning the car by always switching?
What is the probability of winning the car by switching given the player has picked door 1 and the host has opened door 3?
The answer to the first question is , as is shown correctly by the "simple" solutions. But the answer to the second question is now different: the conditional probability the car is behind door 1 or door 2 given the host has opened door 3 (the door on the right) is . This is because Monty's preference for rightmost doors means that he opens door 3 if the car is behind door 1 (which it is originally with probability ) or if the car is behind door 2 (also originally with probability ). For this variation, the two questions yield different answers. This is partially because the assumed condition of the second question (that the host opens door 3) would only occur in this variant with probability . However, as long as the initial probability the car is behind each door is , it is never to the contestant's disadvantage to switch, as the conditional probability of winning by switching is always at least .
In Morgan et al., four university professors published an article in The American Statistician claiming that Savant gave the correct advice but the wrong argument. They believed the question asked for the chance of the car behind door 2 given the player's initial choice of door 1 and the game host opening door 3, and they showed this chance was anything between and 1 depending on the host's decision process given the choice. Only when the decision is completely randomized is the chance .
In an invited comment and in subsequent letters to the editor, Morgan et al were supported by some writers, criticized by others; in each case a response by Morgan et al is published alongside the letter or comment in The American Statistician. In particular, Savant defended herself vigorously. Morgan et al complained in their response to Savant that Savant still had not actually responded to their own main point. Later in their response to Hogbin and Nijdam, they did agree that it was natural to suppose that the host chooses a door to open completely at random when he does have a choice, and hence that the conditional probability of winning by switching (i.e., conditional given the situation the player is in when he has to make his choice) has the same value, , as the unconditional probability of winning by switching (i.e., averaged over all possible situations). This equality was already emphasized by Bell (1992), who suggested that Morgan et als mathematically-involved solution would appeal only to statisticians, whereas the equivalence of the conditional and unconditional solutions in the case of symmetry was intuitively obvious.
There is disagreement in the literature regarding whether Savant's formulation of the problem, as presented in Parade, is asking the first or second question, and whether this difference is significant. Behrends concludes that "One must consider the matter with care to see that both analyses are correct", which is not to say that they are the same. Several critics of the paper by Morgan et al., whose contributions were published along with the original paper, criticized the authors for altering Savant's wording and misinterpreting her intention. One discussant (William Bell) considered it a matter of taste whether one explicitly mentions that (by the standard conditions) which door is opened by the host is independent of whether one should want to switch.
Among the simple solutions, the "combined doors solution" comes closest to a conditional solution, as we saw in the discussion of methods using the concept of odds and Bayes' theorem. It is based on the deeply rooted intuition that revealing information that is already known does not affect probabilities. But, knowing that the host can open one of the two unchosen doors to show a goat does not mean that opening a specific door would not affect the probability that the car is behind the door chosen initially. The point is, though we know in advance that the host will open a door and reveal a goat, we do not know which door he will open. If the host chooses uniformly at random between doors hiding a goat (as is the case in the standard interpretation), this probability indeed remains unchanged, but if the host can choose non-randomly between such doors, then the specific door that the host opens reveals additional information. The host can always open a door revealing a goat and (in the standard interpretation of the problem) the probability that the car is behind the initially chosen door does not change, but it is not because of the former that the latter is true. Solutions based on the assertion that the host's actions cannot affect the probability that the car is behind the initially chosen appear persuasive, but the assertion is simply untrue unless both of the host's two choices are equally likely, if he has a choice. The assertion therefore needs to be justified; without justification being given, the solution is at best incomplete. It can be the case that the answer is correct but the reasoning used to justify it is defective.
Solutions using conditional probability and other solutions
The simple solutions above show that a player with a strategy of switching wins the car with overall probability , i.e., without taking account of which door was opened by the host. In accordance with this, most sources for the topic of probability calculate the conditional probabilities that the car is behind door 1 and door 2 to be and respectively given the contestant initially picks door 1 and the host opens door 3. The solutions in this section consider just those cases in which the player picked door 1 and the host opened door 3.
Refining the simple solution
If we assume that the host opens a door at random, when given a choice, then which door the host opens gives us no information at all as to whether or not the car is behind door 1. In the simple solutions, we have already observed that the probability that the car is behind door 1, the door initially chosen by the player, is initially . Moreover, the host is certainly going to open a (different) door, so opening a door (which door is unspecified) does not change this. must be the average of: the probability that the car is behind door 1, given that the host picked door 2, and the probability of car behind door 1, given the host picked door 3: this is because these are the only two possibilities. However, these two probabilities are the same. Therefore, they are both equal to . This shows that the chance that the car is behind door 1, given that the player initially chose this door and given that the host opened door 3, is , and it follows that the chance that the car is behind door 2, given that the player initially chose door 1 and the host opened door 3, is . The analysis also shows that the overall success rate of , achieved by always switching, cannot be improved, and underlines what already may well have been intuitively obvious: the choice facing the player is that between the door initially chosen, and the other door left closed by the host, the specific numbers on these doors are irrelevant.
Conditional probability by direct calculation
By definition, the conditional probability of winning by switching given the contestant initially picks door 1 and the host opens door 3 is the probability for the event "car is behind door 2 and host opens door 3" divided by the probability for "host opens door 3". These probabilities can be determined referring to the conditional probability table below, or to an equivalent decision tree. The conditional probability of winning by switching is , which is .
The conditional probability table below shows how 300 cases, in all of which the player initially chooses door 1, would be split up, on average, according to the location of the car and the choice of door to open by the host.
Bayes' theorem
Many probability text books and articles in the field of probability theory derive the conditional probability solution through a formal application of Bayes' theoremamong them books by Gill and Henze. Use of the odds form of Bayes' theorem, often called Bayes' rule, makes such a derivation more transparent.
Initially, the car is equally likely to be behind any of the three doors: the odds on door 1, door 2, and door 3 are . This remains the case after the player has chosen door 1, by independence. According to Bayes' rule, the posterior odds on the location of the car, given that the host opens door 3, are equal to the prior odds multiplied by the Bayes factor or likelihood, which is, by definition, the probability of the new piece of information (host opens door 3) under each of the hypotheses considered (location of the car). Now, since the player initially chose door 1, the chance that the host opens door 3 is 50% if the car is behind door 1, 100% if the car is behind door 2, 0% if the car is behind door 3. Thus the Bayes factor consists of the ratios or equivalently , while the prior odds were . Thus, the posterior odds become equal to the Bayes factor . Given that the host opened door 3, the probability that the car is behind door 3 is zero, and it is twice as likely to be behind door 2 than door 1.
Richard Gill analyzes the likelihood for the host to open door 3 as follows. Given that the car is not behind door 1, it is equally likely that it is behind door 2 or 3. Therefore, the chance that the host opens door 3 is 50%. Given that the car is behind door 1, the chance that the host opens door 3 is also 50%, because, when the host has a choice, either choice is equally likely. Therefore, whether or not the car is behind door 1, the chance that the host opens door 3 is 50%. The information "host opens door 3" contributes a Bayes factor or likelihood ratio of , on whether or not the car is behind door 1. Initially, the odds against door 1 hiding the car were . Therefore, the posterior odds against door 1 hiding the car remain the same as the prior odds, .
In words, the information which door is opened by the host (door 2 or door 3?) reveals no information at all about whether or not the car is behind door 1, and this is precisely what is alleged to be intuitively obvious by supporters of simple solutions, or using the idioms of mathematical proofs, "obviously true, by symmetry".
Strategic dominance solution
Going back to Nalebuff, the Monty Hall problem is also much studied in the literature on game theory and decision theory, and also some popular solutions correspond to this point of view. Savant asks for a decision, not a chance. And the chance aspects of how the car is hidden and how an unchosen door is opened are unknown. From this point of view, one has to remember that the player has two opportunities to make choices: first of all, which door to choose initially, and secondly, whether or not to switch. Since he does not know how the car is hidden nor how the host makes choices, he may be able to make use of his first choice opportunity, as it were to neutralize the actions of the team running the quiz show, including the host.
Following Gill, a strategy of contestant involves two actions: the initial choice of a door and the decision to switch (or to stick) which may depend on both the door initially chosen and the door to which the host offers switching. For instance, one contestant's strategy is "choose door 1, then switch to door 2 when offered, and do not switch to door 3 when offered". Twelve such deterministic strategies of the contestant exist.
Elementary comparison of contestant's strategies shows that, for every strategy A, there is another strategy B "pick a door then switch no matter what happens" that dominates it. No matter how the car is hidden and no matter which rule the host uses when he has a choice between two goats, if A wins the car then B also does. For example, strategy A "pick door 1 then always stick with it" is dominated by the strategy B "pick door 2 then always switch after the host reveals a door": A wins when door 1 conceals the car, while B wins when either of the doors 1 or 3 conceals the car.
Similarly, strategy A "pick door 1 then switch to door 2 (if offered), but do not switch to door 3 (if offered)" is dominated by strategy B "pick door 2 then always switch". A wins when door 1 conceals the car and Monty chooses to open door 2 or if door 3 conceals the car. Strategy B wins when either door 1 or door 3 conceals the car, that is, whenever A wins plus the case where door 1 conceals the car and Monty chooses to open door 3.
Dominance is a strong reason to seek for a solution among always-switching strategies, under fairly general assumptions on the environment in which the contestant is making decisions. In particular, if the car is hidden by means of some randomization devicelike tossing symmetric or asymmetric three-sided diethe dominance implies that a strategy maximizing the probability of winning the car will be among three always-switching strategies, namely it will be the strategy that initially picks the least likely door then switches no matter which door to switch is offered by the host.
Strategic dominance links the Monty Hall problem to game theory. In the zero-sum game setting of Gill, discarding the non-switching strategies reduces the game to the following simple variant: the host (or the TV-team) decides on the door to hide the car, and the contestant chooses two doors (i.e., the two doors remaining after the player's first, nominal, choice). The contestant wins (and her opponent loses) if the car is behind one of the two doors she chose.
Solutions by simulation
A simple way to demonstrate that a switching strategy really does win two out of three times with the standard assumptions is to simulate the game with playing cards. Three cards from an ordinary deck are used to represent the three doors; one 'special' card represents the door with the car and two other cards represent the goat doors.
The simulation can be repeated several times to simulate multiple rounds of the game. The player picks one of the three cards, then, looking at the remaining two cards the 'host' discards a goat card. If the card remaining in the host's hand is the car card, this is recorded as a switching win; if the host is holding a goat card, the round is recorded as a staying win. As this experiment is repeated over several rounds, the observed win rate for each strategy is likely to approximate its theoretical win probability, in line with the law of large numbers.
Repeated plays also make it clearer why switching is the better strategy. After the player picks his card, it is already determined whether switching will win the round for the player. If this is not convincing, the simulation can be done with the entire deck. In this variant, the car card goes to the host 51 times out of 52, and stays with the host no matter how many non-car cards are discarded.
Variants
A common variant of the problem, assumed by several academic authors as the canonical problem, does not make the simplifying assumption that the host must uniformly choose the door to open, but instead that he uses some other strategy. The confusion as to which formalization is authoritative has led to considerable acrimony, particularly because this variant makes proofs more involved without altering the optimality of the always-switch strategy for the player. In this variant, the player can have different probabilities of winning depending on the observed choice of the host, but in any case the probability of winning by switching is at least (and can be as high as 1), while the overall probability of winning by switching is still exactly . The variants are sometimes presented in succession in textbooks and articles intended to teach the basics of probability theory and game theory. A considerable number of other generalizations have also been studied.
Other host behaviors
The version of the Monty Hall problem published in Parade in 1990 did not specifically state that the host would always open another door, or always offer a choice to switch, or even never open the door revealing the car. However, Savant made it clear in her second follow-up column that the intended host's behavior could only be what led to the probability she gave as her original answer. "Anything else is a different question." "Virtually all of my critics understood the intended scenario. I personally read nearly three thousand letters (out of the many additional thousands that arrived) and found nearly every one insisting simply that because two options remained (or an equivalent error), the chances were even. Very few raised questions about ambiguity, and the letters actually published in the column were not among those few." The answer follows if the car is placed randomly behind any door, the host must open a door revealing a goat regardless of the player's initial choice and, if two doors are available, chooses which one to open randomly. The table below shows a variety of other possible host behaviors and the impact on the success of switching.
Determining the player's best strategy within a given set of other rules the host must follow is the type of problem studied in game theory. For example, if the host is not required to make the offer to switch the player may suspect the host is malicious and makes the offers more often if the player has initially selected the car. In general, the answer to this sort of question depends on the specific assumptions made about the host's behavior, and might range from "ignore the host completely" to "toss a coin and switch if it comes up heads"; see the last row of the table below.
Morgan et al and Gillman both show a more general solution where the car is (uniformly) randomly placed but the host is not constrained to pick uniformly randomly if the player has initially selected the car, which is how they both interpret the statement of the problem in Parade despite the author's disclaimers. Both changed the wording of the Parade version to emphasize that point when they restated the problem. They consider a scenario where the host chooses between revealing two goats with a preference expressed as a probability q, having a value between 0 and 1. If the host picks randomly q would be and switching wins with probability regardless of which door the host opens. If the player picks door 1 and the host's preference for door 3 is q, then the probability the host opens door 3 and the car is behind door 2 is , while the probability the host opens door 3 and the car is behind door 1 is . These are the only cases where the host opens door 3, so the conditional probability of winning by switching given the host opens door 3 is which simplifies to . Since q can vary between 0 and 1 this conditional probability can vary between and 1. This means even without constraining the host to pick randomly if the player initially selects the car, the player is never worse off switching. However neither source suggests the player knows what the value of q is so the player cannot attribute a probability other than the that Savant assumed was implicit.
N doors
D. L. Ferguson (1975 in a letter to Selvin) suggests an N-door generalization of the original problem in which the host opens p losing doors and then offers the player the opportunity to switch; in this variant switching wins with probability . This probability is always greater than , therefore switching always brings an advantage.
Even if the host opens only a single door (), the player is better off switching in every case. As N grows larger, the advantage decreases and approaches zero.
At the other extreme, if the host opens all losing doors but one (p = N − 2) the advantage increases as N grows large (the probability of winning by switching is , which approaches 1 as N grows very large).
Quantum version
A quantum version of the paradox illustrates some points about the relation between classical or non-quantum information and quantum information, as encoded in the states of quantum mechanical systems. The formulation is loosely based on quantum game theory. The three doors are replaced by a quantum system allowing three alternatives; opening a door and looking behind it is translated as making a particular measurement. The rules can be stated in this language, and once again the choice for the player is to stick with the initial choice, or change to another "orthogonal" option. The latter strategy turns out to double the chances, just as in the classical case. However, if the show host has not randomized the position of the prize in a fully quantum mechanical way, the player can do even better, and can sometimes even win the prize with certainty.
History
The earliest of several probability puzzles related to the Monty Hall problem is Bertrand's box paradox, posed by Joseph Bertrand in 1889 in his Calcul des probabilités. In this puzzle, there are three boxes: a box containing two gold coins, a box with two silver coins, and a box with one of each. After choosing a box at random and withdrawing one coin at random that happens to be a gold coin, the question is what is the probability that the other coin is gold. As in the Monty Hall problem, the intuitive answer is , but the probability is actually .
The Three Prisoners problem, published in Martin Gardner's Mathematical Games column in Scientific American in 1959 is equivalent to the Monty Hall problem. This problem involves three condemned prisoners, a random one of whom has been secretly chosen to be pardoned. One of the prisoners begs the warden to tell him the name of one of the others to be executed, arguing that this reveals no information about his own fate but increases his chances of being pardoned from to . The warden obliges, (secretly) flipping a coin to decide which name to provide if the prisoner who is asking is the one being pardoned. The question is whether knowing the warden's answer changes the prisoner's chances of being pardoned. This problem is equivalent to the Monty Hall problem; the prisoner asking the question still has a chance of being pardoned but his unnamed colleague has a chance.
Steve Selvin posed the Monty Hall problem in a pair of letters to The American Statistician in 1975. The first letter presented the problem in a version close to its presentation in Parade 15 years later. The second appears to be the first use of the term "Monty Hall problem". The problem is actually an extrapolation from the game show. Monty Hall did open a wrong door to build excitement, but offered a known lesser prizesuch as $100 cashrather than a choice to switch doors. As Monty Hall wrote to Selvin:
A version of the problem very similar to the one that appeared three years later in Parade was published in 1987 in the Puzzles section of The Journal of Economic Perspectives. Nalebuff, as later writers in mathematical economics, sees the problem as a simple and amusing exercise in game theory.
"The Monty Hall Trap", Phillip Martin's 1989 article in Bridge Today, presented Selvin's problem as an example of what Martin calls the probability trap of treating non-random information as if it were random, and relates this to concepts in the game of bridge.
A restated version of Selvin's problem appeared in Marilyn vos Savant's Ask Marilyn question-and-answer column of Parade in September 1990. Though Savant gave the correct answer that switching would win two-thirds of the time, she estimates the magazine received 10,000 letters including close to 1,000 signed by PhDs, many on letterheads of mathematics and science departments, declaring that her solution was wrong. Due to the overwhelming response, Parade published an unprecedented four columns on the problem. As a result of the publicity the problem earned the alternative name "Marilyn and the Goats".
In November 1990, an equally contentious discussion of Savant's article took place in Cecil Adams's column "The Straight Dope". Adams initially answered, incorrectly, that the chances for the two remaining doors must each be one in two. After a reader wrote in to correct the mathematics of Adams's analysis, Adams agreed that mathematically he had been wrong. "You pick door #1. Now you're offered this choice: open door #1, or open door #2 and door #3. In the latter case you keep the prize if it's behind either door. You'd rather have a two-in-three shot at the prize than one-in-three, wouldn't you? If you think about it, the original problem offers you basically the same choice. Monty is saying in effect: you can keep your one door or you can have the other two doors, one of which (a non-prize door) I'll open for you." Adams did say the Parade version left critical constraints unstated, and without those constraints, the chances of winning by switching were not necessarily two out of three (e.g., it was not reasonable to assume the host always opens a door). Numerous readers, however, wrote in to claim that Adams had been "right the first time" and that the correct chances were one in two.
The Parade column and its response received considerable attention in the press, including a front-page story in The New York Times in which Monty Hall himself was interviewed. Hall understood the problem, giving the reporter a demonstration with car keys and explaining how actual game play on Let's Make a Deal differed from the rules of the puzzle. In the article, Hall pointed out that because he had control over the way the game progressed, playing on the psychology of the contestant, the theoretical solution did not apply to the show's actual gameplay. He said he was not surprised at the experts' insistence that the probability was 1 out of 2. "That's the same assumption contestants would make on the show after I showed them there was nothing behind one door," he said. "They'd think the odds on their door had now gone up to 1 in 2, so they hated to give up the door no matter how much money I offered. By opening that door we were applying pressure. We called it the Henry James treatment. It was 'The Turn of the Screw'." Hall clarified that as a game show host he did not have to follow the rules of the puzzle in the Savant column and did not always have to allow a person the opportunity to switch (e.g., he might open their door immediately if it was a losing door, might offer them money to not switch from a losing door to a winning door, or might allow them the opportunity to switch only if they had a winning door). "If the host is required to open a door all the time and offer you a switch, then you should take the switch," he said. "But if he has the choice whether to allow a switch or not, beware. Caveat emptor. It all depends on his mood."
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6028064 | https://en.wikipedia.org/wiki/Lymphatic%20filariasis | Lymphatic filariasis | Lymphatic filariasis is a human disease caused by parasitic worms known as filarial worms. Usually acquired in childhood, it is a leading cause of permanent disability worldwide, impacting over a hundred million people and manifesting itself in a variety of severe clinical pathologies While most cases have no symptoms, some people develop a syndrome called elephantiasis, which is marked by severe swelling in the arms, legs, breasts, or genitals. The skin may become thicker as well, and the condition may become painful. Affected people are often unable to work and are often shunned or rejected by others because of their disfigurement and disability.
It is the first of the mosquito-borne diseases to have been identified. The worms are spread by the bites of infected mosquitoes. Three types of worms are known to cause the disease: Wuchereria bancrofti, Brugia malayi, and Brugia timori, with Wuchereria bancrofti being the most common. These worms damage the lymphatic system by nesting within the lymphatic vessels and disrupting the system's normal function. Worms can survive within the human body for up to 8 years, all while reproducing millions of larvae which circulate through the blood. The disease is diagnosed by microscopic examination of blood collected during the night. The blood is typically examined as a smear after being stained with Giemsa stain. Testing the blood for antibodies against the disease may also permit diagnosis. Other roundworms from the same family are responsible for river blindness.
Prevention can be achieved by treating entire groups affected by the disease, known as mass deworming. This is done every year for about six years, to rid a population of the disease entirely. Medications usually include a combination of two or more anthelmintic agents: albendazole, ivermectin, and diethylcarbamazine. Efforts to prevent mosquito bites are also recommended, including reducing the number of mosquitoes and promoting the use of bed nets.
As of 2022, about 40 million people were infected, and about 863 million people were at risk of the disease in 47 countries. It is most common in tropical Africa and Asia. Lymphatic filariasis is classified as a neglected tropical disease and one of the four main worm infections. The impact of the disease results in economic losses of billions of US dollars a year.
Signs and symptoms
Most people infected with the worms that cause lymphatic filariasis never develop symptoms; though some have damage to lymph vessels that can be detected by medical ultrasound. Months to years after the initial infection, the worms die, triggering an immune response that manifests with repeated episodes of fever and painful swelling over the nearest lymph nodes (typically those along the groin). In areas with endemic lymphatic filariasis, people are typically infected in childhood, and symptoms begin in adolescence.
A subset of those affected have continued damage to their lymph vessels. Dysfunctional vessels fail to recirculate lymph fluid, which can pool (called lymphodema) in the nearest extremity – generally the arm, leg, breast, or scrotum. Loss of lymph function (which transports immune cells) results in various repeated infections in the area. Repeated cycles of infection, inflammation, and lymph vessel damage over several years cause the affected extremity to swell to an extremely large size. The surrounding skin thickens, becoming dry, discolored, and dotted with wartlike lumps that contain tortuous loops of lymph vessels.
Even those without lymph damage can sometimes develop an allergic reaction to the worm larvae in the capillaries of the lung, called tropical pulmonary eosinophilia. These people develop nocturnal coughing, fatigue, and weight loss. Over time, this can damage the lungs, resulting in restrictive lung disease (decreased lung capacity).
Causes
Lymphatic filariasis is caused by infection with three different nematode worms: Wuchereria bancrofti (causes 90% of cases), Brugia malayi, and Brugia timori. The three worms are transmitted by the bite of an infected mosquito – largely of genera Aedes, Anopheles, Culex, or Mansonia. When the mosquito bites, infectious nematode larvae are dropped onto the skin. They crawl into the bite wound, through the subcutaneous tissue, and into nearby lymph vessels. There, they develop into adults over about a year, with adult females up to long, and males up to half that length. Adult females and males mate, prompting the female to begin releasing a constant stream of larvae called "microfilariae" – more than 10,000 microfilariae each day for the adult's remaining lifespan of around five to eight years. Microfilariae typically circulate in the blood stream at night; during the day they collect in the capillaries of the lungs.
A mosquito that feeds on an infected person can take up microfilariae along with its blood meal. Inside the mosquito, the microfilariae pierce the stomach wall and crawl to the flight muscles, where they mature over 10 to 20 days into their human-infectious form. They then crawl to the mosquito's mouth to be deposited at its next bite, continuing the lifecycle.
The disease itself is a result of a complex interplay between several factors: the worm, the endosymbiotic Wolbachia bacteria within the worm, the host's immune response, and the numerous opportunistic infections and disorders that arise. The adult worms live in the human lymphatic system and obstruct the flow of lymph throughout the body; this results in chronic lymphedema, most often noted in the lower torso (typically in the legs and genitals). These worms can survive within the human body for up to 8 years, all while reproducing millions of larvae which circulate through the blood.
Diagnosis
The preferred method for diagnosing lymphatic filariasis is by finding the microfilariae via microscopic examination of the blood. The blood sample is typically in the form of a thick smear, stained with Giemsa stain. Technicians analyzing the blood smear must be able to distinguish between W. bancrofti and other parasites potentially present. A blood smear is a simple and fairly accurate diagnostic tool, provided the blood sample is taken when the microfilariae are in the peripheral circulation. Because the microfilariae only circulate in the blood at night, the blood specimen must be collected at night.
It is often difficult or impossible to detect the causative organism in the peripheral blood, even in advanced cases. In such cases, testing the blood serum for antibodies against the disease may also be used. A polymerase chain reaction test can also be performed to detect a minute fraction, as little as 1 pg, of filarial DNA. Dead, calcified worms can be detected by X-ray examinations. Ultrasonography can also be used to detect the movements and noises caused by the movement of adult worms.
Differential diagnosis
Lymphatic filariasis may be confused with podoconiosis (also known as nonfilarial elephantiasis), a non-infectious disease caused by exposure of bare feet to irritant alkaline clay soils. Podoconiosis however typically affects the legs bilaterally, while filariasis is generally unilateral. Also, podoconiosis very rarely affects the groin while filariasis frequently involves the groin. Geographical location may also help to distinguish between these two diseases: podoconiosis is typically found in higher altitude areas with high seasonal rainfall, while filariasis is common in low-lying areas where mosquitos are prevalent.
Prevention
Protecting against mosquito bites in endemic regions is crucial to the prevention of lymphatic filariasis. Insect repellents and mosquito nets (especially when treated with an insecticide such as deltamethrin or permethrin) have been demonstrated to reduce the transmission of lymphatic filariasis. In addition residual spraying and personal protective equipment are known ways to control vectors.
Worldwide eradication of lymphatic filariasis is the definitive goal. This is considered to be achievable since the disease has no known animal reservoir. The World Health Organization (WHO) is coordinating the global effort to eradicate filariasis. The mainstay of this program is mass deworming of entire populations of people who are at risk with antifilarial drugs. The specific treatment depends on the co-endemicity of lymphatic filariasis with other filarial diseases. The WHO's annual MDA guidelines are listed below.
For areas co-endemic with loiasis 400 mg of albendazole should be administered
for countries co-endemic with onchocerciasis, 200 mcg/kg of Ivermectin should be administered with 400 mg of albendazole
in countries without onchocerciasis 6 mg/kg of diethylcarbamazine citrate (DEC) and 400 mg of albendazole should be used
in countries without onchocerciasis and the IDA Guidelines are met 200 mcg/kg of ivermectin should be used with 6 mg/kg of diethylcarbamazine and 400 mg of albendazole.
Because the parasite requires a human host to reproduce, consistent treatment of at-risk populations (annually for a duration of four to six years) is expected to break the cycle of transmission and cause the extinction of the causative organisms.
In 2011, Sri Lanka was certified by the WHO as having eradicated lymphatic filariasis. In July 2017, the WHO announced that the disease had been eliminated in Tonga. Elimination of the disease has also occurred in Cambodia, China, Cook Islands, Egypt, Kiribati, Maldives, Marshall Islands, Niue, Palau, South Korea, Thailand, Vanuatu, Vietnam, and Wallis and Futuna. In 2020, the WHO announced 2030 targets for this program of eliminating lymphatic filariasis in 80% of endemic countries.
A vaccine is not yet available, but in 2013, the University of Illinois College of Medicine was reporting 95% efficacy in testing against B. malayi in mice.
Treatment
Treatment of lymphatic filariasis depends in part on the geographic location of the area of the world in which the disease was acquired but often involves the combination of two or more anthelmintic agents: albendazole, ivermectin, and diethylcarbamazine. In sub-Saharan Africa, the disease is usually treated with albendazole and ivermectin, whereas in the Western Pacific region, all three anthelmintic agents are used. While diethylcarbamazine in combination with albendazole is often used, it isn't as region-specific as the other combinations. Albendazole monotherapy is used in regions with endemic loiasis in combination with integrated vector control and was the treatment modality for 11.2 million people in 2022.
Wolbachia are endosymbiotic bacteria that live inside the gut of the nematode worms responsible for lymphatic filariasis, and which provide nutrients necessary for their survival. Doxycycline kills these bacteria, which in turn prevents the maturation of microfilariae into adults. It also shortens the lifespan of the adult worms, causing them to die within 1 to 2 years instead of their normal 10 to 14-year lifespan. Doxycycline is effective in treating lymphatic filariasis. Limitations of this antibiotic protocol include that it requires 4 to 6 weeks of treatment rather than the single dose of the anthelmintic agents, that doxycycline should not be used in young children and pregnant women, and that it is phototoxic.
Albendazole is classified as an antihelmintics, which specifically works to kill worms. The drug stops the worms from absorbing glucose, evidently leading to starvation and death from fatigue. The effects of albendazole alone have varying results, however, in combination with DEC drugs it has been found more effective. Ivermectin is administered with albendazole, and works by binding to the nerve cells of the parasites, subsequently making them permeable to chloride. This leads to death by paralysis. Ivermectin, however, has been found to only kill the parasites in their early stages of life, and cannot kill an adult, live worm. Therefore, this drug is usually combined with DEC to kill both the microfilariae and the adult worms.
Surgical treatment may be helpful in cases of scrotal elephantiasis and hydrocele. However, surgery is generally ineffective at correcting elephantiasis of the limbs. Acute inflammatory responses due to lymphedema, and hydrocele can be reduced or prevented by practicing good hygiene, skin care, exercise and elevation of infected limbs.
Epidemiology
Lymphatic filariasis occurs in tropical and subtropical regions of Africa, Asia, Central America, the Caribbean, South America, and certain Pacific Island nations. Elephantiasis caused by lymphatic filariasis is one of the most common causes of permanent disability in the world. As of 2018, 51 million people were infected with lymphatic filariasis and at least 863 million people in 50 countries were living in areas that require preventive chemotherapy to stop the spread of infection. By 2022, the prevalence had declined to somewhere around 40 million and the disease remains endemic in 47 countries. These improvements are a direct result of the WHO's Global Programme to Eliminate Lymphatic Filariasis. Since implementation, 740 million people no longer require preventative chemotherapy to treat the disease.
W. bancrofti is responsible for 90% of lymphatic filariasis. Brugia malayi causes most of the remainder of the cases, while Brugia timori is a rare cause. W. bancrofti largely affects areas across the broad equatorial belt (Africa, the Nile Delta, Turkey, India, the East Indies, Southeast Asia, Philippines, Oceanic Islands, and parts of South America). Since lymphatic filariasis requires multiple mosquito bites over several months to years to spread of infection due to tourism is low. The mosquito vectors of W. bancrofti have a preference for human blood; humans are the only animals naturally infected with W. bancrofti. No reservoir host is known. Lymphatic Filariasis is extremely uncommon in the United States, with only one reported case found in South Carolina in the early 1900s.
In South America, four endemic countries have been working to beset lymphatic filariasis, consisting of Brazil, the Dominican Republic, Guyana, and Haiti. In Latin America, the spread of lymphatic filariasis is through W. bancrofti, the only anthropods within the region, Culex quinquefasciatus. The exponential rate of development within the Americas is being combated through the development of an MDA program. MDA program, a 3-step drug administering program, has led to a 67% decrease in the need for the drug program. Brazil targeted the rising endemic by administering DEC drugs through an MDA program to the communities hit hardest by the disease. By providing these drugs annually, as well as offering post-care, through showing family members how to treat the disease, creating connections for jobs, as well as providing a social network to incorporate patients into society, Brazil has made the most effort to provide care. Dominican Republic has administered 5 rounds of DEC drugs annually for five years, spanning from 2002-2007. After the initial drastic action, the Dominican then administered another three rounds of MDA. Guyana also used DEC drugs to focus on preventing the spread of the disease, using a DEC-fortified salt from 2003-2007 and ultimately switching to MDA with DEC from 2014 to the present. Targeting patient education and access to treatment. Haiti then focused on the disease by implementing DEC drug in 2002. It reached full geographical coverage by 2012, subsequently in 2014 about 20 communities had eradicated the need for MDA.
In communities where lymphatic filariasis is endemic, as many as 10% of women can be affected by swollen limbs, and 50% of men can develop mutilating genital symptoms.
History
There is evidence of Lymphatic filariasis cases dating back 4000 years. The ancient Vedic text, the Rig Veda, composed around 1500 BC–1200 BC, makes a possible reference to elephantiasis. The 50th hymn of the 7th book of the Rigveda calls on the gods Mitra, Varuna, and Agni for protection against "that which nests inside and swells". The author of the hymn implores the deities to not let the worm wound his foot. The disease is described as causing eruptions to appear on the ankles and the knees. Artifacts from ancient Egypt (2000 BC) and the Nok civilization in West Africa (500 BC) show possible elephantiasis symptoms. The first clear reference to the disease occurs in ancient Greek literature, wherein scholars differentiated the often similar symptoms of lymphatic filariasis from those of leprosy, describing leprosy as elephantiasis graecorum and lymphatic filariasis as elephantiasis arabum.
The first documentation of symptoms occurred in the 16th century, when Jan Huyghen van Linschoten wrote about the disease during the exploration of Goa. Similar symptoms were reported by subsequent explorers in areas of Asia and Africa, though an understanding of the disease did not begin to develop until centuries later.
The causative agents were first identified in the late 19th century. In 1866, Timothy Lewis, building on the work of and Otto Henry Wucherer, made the connection between microfilariae and elephantiasis, establishing the course of research that would ultimately explain the disease. In 1876, Joseph Bancroft discovered the adult form of the worm. In 1877, the lifecycle involving an arthropod vector was theorized by Patrick Manson, who proceeded to demonstrate the presence of the worms in mosquitoes. Manson incorrectly hypothesized that the disease was transmitted through skin contact with water in which the mosquitoes had laid eggs. In 1900, George Carmichael Low determined the actual transmission method by discovering the presence of the worm in the proboscis of the mosquito vector.
Research directions
Researchers at the University of Illinois at Chicago (UIC) have developed a novel vaccine for the prevention of lymphatic filariasis. This vaccine has been shown to elicit strong, protective immune responses in mouse models of lymphatic filariasis infection. The immune response elicited by this vaccine has been demonstrated to be protective against both W. bancrofti and B. malayi infection in the mouse model and may prove useful in the human.
On 20 September 2007, geneticists published the first draft of the complete genome (genetic content) of Brugia malayi, one of the roundworms that causes lymphatic filariasis. This project had been started in 1994 and by 2000, 80% of the genome had been determined. Determining the content of the genes might lead to the development of new drugs and vaccines.
Veterinary disease
Onchocerca ochengi causes lymphatic filariasis in cattle.
| Biology and health sciences | Helminthic diseases and infestations | Health |
1789646 | https://en.wikipedia.org/wiki/Mathematics%20of%20general%20relativity | Mathematics of general relativity | When studying and formulating Albert Einstein's theory of general relativity, various mathematical structures and techniques are utilized. The main tools used in this geometrical theory of gravitation are tensor fields defined on a Lorentzian manifold representing spacetime. This article is a general description of the mathematics of general relativity.
Note: General relativity articles using tensors will use the abstract index notation.
Tensors
The principle of general covariance was one of the central principles in the development of general relativity. It states that the laws of physics should take the same mathematical form in all reference frames. The term 'general covariance' was used in the early formulation of general relativity, but the principle is now often referred to as 'diffeomorphism covariance'.
Diffeomorphism covariance is not the defining feature of general relativity,[1] and controversies remain regarding its present status in general relativity. However, the invariance property of physical laws implied in the principle, coupled with the fact that the theory is essentially geometrical in character (making use of non-Euclidean geometries), suggested that general relativity be formulated using the language of tensors. This will be discussed further below.
Spacetime as a manifold
Most modern approaches to mathematical general relativity begin with the concept of a manifold. More precisely, the basic physical construct representing a curved is modelled by a four-dimensional, smooth, connected, Lorentzian manifold. Other physical descriptors are represented by various tensors, discussed below.
The rationale for choosing a manifold as the fundamental mathematical structure is to reflect desirable physical properties. For example, in the theory of manifolds, each point is contained in a (by no means unique) coordinate chart, and this chart can be thought of as representing the 'local spacetime' around the observer (represented by the point). The principle of local Lorentz covariance, which states that the laws of special relativity hold locally about each point of spacetime, lends further support to the choice of a manifold structure for representing spacetime, as locally around a point on a general manifold, the region 'looks like', or approximates very closely Minkowski space (flat spacetime).
The idea of coordinate charts as 'local observers who can perform measurements in their vicinity' also makes good physical sense, as this is how one actually collects physical data locally. For cosmological problems, a coordinate chart may be quite large.
Local versus global structure
An important distinction in physics is the difference between local and global structures. Measurements in physics are performed in a relatively small region of spacetime and this is one reason for studying the local structure of spacetime in general relativity, whereas determining the global spacetime structure is important, especially in cosmological problems.
An important problem in general relativity is to tell when two spacetimes are 'the same', at least locally. This problem has its roots in manifold theory where determining if two Riemannian manifolds of the same dimension are locally isometric ('locally the same'). This latter problem has been solved and its adaptation for general relativity is called the Cartan–Karlhede algorithm.
Tensors in general relativity
One of the profound consequences of relativity theory was the abolition of privileged reference frames. The description of physical phenomena should not depend upon who does the measuring one reference frame should be as good as any other. Special relativity demonstrated that no inertial reference frame was preferential to any other inertial reference frame, but preferred inertial reference frames over noninertial reference frames. General relativity eliminated preference for inertial reference frames by showing that there is no preferred reference frame (inertial or not) for describing nature.
Any observer can make measurements and the precise numerical quantities obtained only depend on the coordinate system used. This suggested a way of formulating relativity using 'invariant structures', those that are independent of the coordinate system (represented by the observer) used, yet still have an independent existence. The most suitable mathematical structure seemed to be a tensor. For example, when measuring the electric and magnetic fields produced by an accelerating charge, the values of the fields will depend on the coordinate system used, but the fields are regarded as having an independent existence, this independence represented by the electromagnetic tensor .
Mathematically, tensors are generalised linear operators multilinear maps. As such, the ideas of linear algebra are employed to study tensors.
At each point of a manifold, the tangent and cotangent spaces to the manifold at that point may be constructed. Vectors (sometimes referred to as contravariant vectors) are defined as elements of the tangent space and covectors (sometimes termed covariant vectors, but more commonly dual vectors or one-forms) are elements of the cotangent space.
At , these two vector spaces may be used to construct type tensors, which are real-valued multilinear maps acting on the direct sum of copies of the cotangent space with copies of the tangent space. The set of all such multilinear maps forms a vector space, called the tensor product space of type at and denoted by If the tangent space is n-dimensional, it can be shown that
In the general relativity literature, it is conventional to use the component syntax for tensors.
A type tensor may be written as
where is a basis for the i-th tangent space and a basis for the j-th cotangent space.
As spacetime is assumed to be four-dimensional, each index on a tensor can be one of four values. Hence, the total number of elements a tensor possesses equals 4R, where R is the count of the number of covariant and contravariant indices on the tensor, (a number called the rank of the tensor).
Symmetric and antisymmetric tensors
Some physical quantities are represented by tensors not all of whose components are independent. Important examples of such tensors include symmetric and antisymmetric tensors. Antisymmetric tensors are commonly used to represent rotations (for example, the vorticity tensor).
Although a generic rank R tensor in 4 dimensions has 4R components, constraints on the tensor such as symmetry or antisymmetry serve to reduce the number of distinct components. For example, a symmetric rank two tensor satisfies and possesses 10 independent components, whereas an antisymmetric (skew-symmetric) rank two tensor satisfies and has 6 independent components. For ranks greater than two, the symmetric or antisymmetric index pairs must be explicitly identified.
Antisymmetric tensors of rank 2 play important roles in relativity theory. The set of all such tensors - often called bivectors forms a vector space of dimension 6, sometimes called bivector space.
The metric tensor
The metric tensor is a central object in general relativity that describes the local geometry of spacetime (as a result of solving the Einstein field equations). Using the weak-field approximation, the metric tensor can also be thought of as representing the 'gravitational potential'. The metric tensor is often just called 'the metric'.
The metric is a symmetric tensor and is an important mathematical tool. As well as being used to raise and lower tensor indices, it also generates the connections which are used to construct the geodesic equations of motion and the Riemann curvature tensor.
A convenient means of expressing the metric tensor in combination with the incremental intervals of coordinate distance that it relates to is through the line element:
This way of expressing the metric was used by the pioneers of differential geometry. While some relativists consider the notation to be somewhat old-fashioned, many readily switch between this and the alternative notation:
The metric tensor is commonly written as a 4×4 matrix. This matrix is symmetric and thus has 10 independent components.
Invariants
One of the central features of GR is the idea of invariance of physical laws. This invariance can be described in many ways, for example, in terms of local Lorentz covariance, the general principle of relativity or diffeomorphism covariance.
A more explicit description can be given using tensors. The crucial feature of tensors used in this approach is the fact that (once a metric is given) the operation of contracting a tensor of rank R over all R indices gives a number an invariant that is independent of the coordinate chart one uses to perform the contraction. Physically, this means that if the invariant is calculated by any two observers, they will get the same number, thus suggesting that the invariant has some independent significance. Some important invariants in relativity include:
The Ricci scalar:
The Kretschmann scalar:
Other examples of invariants in relativity include the electromagnetic invariants, and various other curvature invariants, some of the latter finding application in the study of gravitational entropy and the Weyl curvature hypothesis.
Tensor classifications
The classification of tensors is a purely mathematical problem. In GR, however, certain tensors that have a physical interpretation can be classified with the different forms of the tensor usually corresponding to some physics. Examples of tensor classifications useful in general relativity include the Segre classification of the energy–momentum tensor and the Petrov classification of the Weyl tensor. There are various methods of classifying these tensors, some of which use tensor invariants.
Tensor fields in general relativity
Tensor fields on a manifold are maps which attach a tensor to each point of the manifold. This notion can be made more precise by introducing the idea of a fibre bundle, which in the present context means to collect together all the tensors at all points of the manifold, thus 'bundling' them all into one grand object called the tensor bundle. A tensor field is then defined as a map from the manifold to the tensor bundle, each point being associated with a tensor at .
The notion of a tensor field is of major importance in GR. For example, the geometry around a star is described by a metric tensor at each point, so at each point of the spacetime the value of the metric should be given to solve for the paths of material particles. Another example is the values of the electric and magnetic fields (given by the electromagnetic field tensor) and the metric at each point around a charged black hole to determine the motion of a charged particle in such a field.
Vector fields are contravariant rank one tensor fields. Important vector fields in relativity include the four-velocity, , which is the coordinate distance travelled per unit of proper time, the four-acceleration and the four-current describing the charge and current densities. Other physically important tensor fields in relativity include the following:
The stress–energy tensor , a symmetric rank-two tensor.
The electromagnetic field tensor , a rank-two antisymmetric tensor.
Although the word 'tensor' refers to an object at a point, it is common practice to refer to tensor fields on a spacetime (or a region of it) as just 'tensors'.
At each point of a spacetime on which a metric is defined, the metric can be reduced to the Minkowski form using Sylvester's law of inertia.
Tensorial derivatives
Before the advent of general relativity, changes in physical processes were generally described by partial derivatives, for example, in describing changes in electromagnetic fields (see Maxwell's equations). Even in special relativity, the partial derivative is still sufficient to describe such changes. However, in general relativity, it is found that derivatives which are also tensors must be used. The derivatives have some common features including that they are derivatives along integral curves of vector fields.
The problem in defining derivatives on manifolds that are not flat is that there is no natural way to compare vectors at different points. An extra structure on a general manifold is required to define derivatives. Below are described two important derivatives that can be defined by imposing an additional structure on the manifold in each case.
Affine connections
The curvature of a spacetime can be characterised by taking a vector at some point and parallel transporting it along a curve on the spacetime. An affine connection is a rule which describes how to legitimately move a vector along a curve on the manifold without changing its direction.
By definition, an affine connection is a bilinear map , where is a space of all vector fields on the spacetime. This bilinear map can be described in terms of a set of connection coefficients (also known as Christoffel symbols) specifying what happens to components of basis vectors under infinitesimal parallel transport:
Despite their appearance, the connection coefficients are not the components of a tensor.
Generally speaking, there are independent connection coefficients at each point of spacetime. The connection is called symmetric or torsion-free, if . A symmetric connection has at most unique coefficients.
For any curve and two points and on this curve, an affine connection gives rise to a map of vectors in the tangent space at into vectors in the tangent space at :
and can be computed component-wise by solving the differential equation
where is the vector tangent to the curve at the point .
An important affine connection in general relativity is the Levi-Civita connection, which is a symmetric connection obtained from parallel transporting a tangent vector along a curve whilst keeping the inner product of that vector constant along the curve. The resulting connection coefficients (Christoffel symbols) can be calculated directly from the metric. For this reason, this type of connection is often called a metric connection.
The covariant derivative
Let be a point, a vector located at , and a vector field. The idea of differentiating at along the direction of in a physically meaningful way can be made sense of by choosing an affine connection and a parameterized smooth curve such that and . The formula
for a covariant derivative of along associated with connection turns out to give curve-independent results and can be used as a "physical definition" of a covariant derivative.
It can be expressed using connection coefficients:
The expression in brackets, called a covariant derivative of (with respect to the connection) and denoted by , is more often used in calculations:
A covariant derivative of can thus be viewed as a differential operator acting on a vector field sending it to a type tensor (increasing the covariant index by 1) and can be generalised to act on type tensor fields sending them to type tensor fields. Notions of parallel transport can then be defined similarly as for the case of vector fields. By definition, a covariant derivative of a scalar field is equal to the regular derivative of the field.
In the literature, there are three common methods of denoting covariant differentiation:
Many standard properties of regular partial derivatives also apply to covariant derivatives:
In general relativity, one usually refers to "the" covariant derivative, which is the one associated with Levi-Civita affine connection. By definition, Levi-Civita connection preserves the metric under parallel transport, therefore, the covariant derivative gives zero when acting on a metric tensor (as well as its inverse). It means that we can take the (inverse) metric tensor in and out of the derivative and use it to raise and lower indices:
The Lie derivative
Another important tensorial derivative is the Lie derivative. Unlike the covariant derivative, the Lie derivative is independent of the metric, although in general relativity one usually uses an expression that seemingly depends on the metric through the affine connection. Whereas the covariant derivative required an affine connection to allow comparison between vectors at different points, the Lie derivative uses a congruence from a vector field to achieve the same purpose. The idea of Lie dragging a function along a congruence leads to a definition of the Lie derivative, where the dragged function is compared with the value of the original function at a given point. The Lie derivative can be defined for type tensor fields and in this respect can be viewed as a map that sends a type to a type tensor.
The Lie derivative is usually denoted by , where is the vector field along whose congruence the Lie derivative is taken.
The Lie derivative of any tensor along a vector field can be expressed through the covariant derivatives of that tensor and vector field. The Lie derivative of a scalar is just the directional derivative:
Higher rank objects pick up additional terms when the Lie derivative is taken. For example, the Lie derivative of a type tensor is
More generally,
In fact in the above expression, one can replace the covariant derivative with any torsion free connection or locally, with the coordinate dependent derivative , showing that the Lie derivative is independent of the metric. The covariant derivative is convenient however because it commutes with raising and lowering indices.
One of the main uses of the Lie derivative in general relativity is in the study of spacetime symmetries where tensors or other geometrical objects are preserved. In particular, Killing symmetry (symmetry of the metric tensor under Lie dragging) occurs very often in the study of spacetimes. Using the formula above, we can write down the condition that must be satisfied for a vector field to generate a Killing symmetry:
The Riemann curvature tensor
A crucial feature of general relativity is the concept of a curved manifold. A useful way of measuring the curvature of a manifold is with an object called the Riemann (curvature) tensor.
This tensor measures curvature by use of an affine connection by considering the effect of parallel transporting a vector between two points along two curves. The discrepancy between the results of these two parallel transport routes is essentially quantified by the Riemann tensor.
This property of the Riemann tensor can be used to describe how initially parallel geodesics diverge. This is expressed by the equation of geodesic deviation and means that the tidal forces experienced in a gravitational field are a result of the curvature of spacetime.
Using the above procedure, the Riemann tensor is defined as a type tensor and when fully written out explicitly contains the Christoffel symbols and their first partial derivatives. The Riemann tensor has 20 independent components. The vanishing of all these components over a region indicates that the spacetime is flat in that region. From the viewpoint of geodesic deviation, this means that initially parallel geodesics in that region of spacetime will stay parallel.
The Riemann tensor has a number of properties sometimes referred to as the symmetries of the Riemann tensor. Of particular relevance to general relativity are the algebraic and differential Bianchi identities.
The connection and curvature of any Riemannian manifold are closely related, the theory of holonomy groups, which are formed by taking linear maps defined by parallel transport around curves on the manifold, providing a description of this relationship.
What the Riemann tensor allows us to do is tell, mathematically, whether a space is flat or, if curved, how much curvature takes place in any given region. In order to derive the Riemann curvature tensor we must first recall the definition of the covariant derivative of a tensor with one and two indices;
For the formation of the Riemann tensor, the covariant derivative is taken twice with the respect to a tensor of rank one. The equation is set up as follows;
Similarly we have:
Subtracting the two equations, swapping dummy indices and using the symmetry of Christoffel symbols leaves:
or
Finally the Riemann curvature tensor is written as
You can contract indices to make the tensor covariant simply by multiplying by the metric, which will be useful when working with Einstein's field equations,
and by further decomposition,
This tensor is called the Ricci tensor which can also be derived by setting and in the Riemann tensor to the same indice and summing over them. Then the curvature scalar can be found by going one step further,
So now we have 3 different objects,
the Riemann curvature tensor: or
the Ricci tensor:
the scalar curvature:
all of which are useful in calculating solutions to Einstein's field equations.
The energy–momentum tensor
The sources of any gravitational field (matter and energy) is represented in relativity by a type symmetric tensor called the energy–momentum tensor. It is closely related to the Ricci tensor. Being a second rank tensor in four dimensions, the energy–momentum tensor may be viewed as a 4 by 4 matrix. The various admissible matrix types, called Jordan forms cannot all occur, as the energy conditions that the energy–momentum tensor is forced to satisfy rule out certain forms.
Energy conservation
In special and general relativity, there is a local law for the conservation of energy–momentum. It can be succinctly expressed by the tensor equation:
This illustrates the rule of thumb that 'partial derivatives go to covariant derivatives'.
The Einstein field equations
The Einstein field equations (EFE) are the core of general relativity theory. The EFE describe how mass and energy (as represented in the stress–energy tensor) are related to the curvature of space-time (as represented in the Einstein tensor). In abstract index notation, the EFE reads as follows:
where is the Einstein tensor, is the cosmological constant, is the metric tensor, is the speed of light in vacuum and is the gravitational constant, which comes from Newton's law of universal gravitation.
The solutions of the EFE are metric tensors. The EFE, being non-linear differential equations for the metric, are often difficult to solve. There are a number of strategies used to solve them. For example, one strategy is to start with an ansatz (or an educated guess) of the final metric, and refine it until it is specific enough to support a coordinate system but still general enough to yield a set of simultaneous differential equations with unknowns that can be solved for. Metric tensors resulting from cases where the resultant differential equations can be solved exactly for a physically reasonable distribution of energy–momentum are called exact solutions. Examples of important exact solutions include the Schwarzschild solution and the Friedman-Lemaître-Robertson–Walker solution.
The EIH approximation plus other references (e.g. Geroch and Jang, 1975 - 'Motion of a body in general relativity', JMP, Vol. 16 Issue 1).
The geodesic equations
Once the EFE are solved to obtain a metric, it remains to determine the motion of inertial objects in the spacetime. In general relativity, it is assumed that inertial motion occurs along timelike and null geodesics of spacetime as parameterized by proper time. Geodesics are curves that parallel transport their own tangent vector ; i.e., . This condition, the geodesic equation, can be written in terms of a coordinate system with the tangent vector :
where denotes the derivative by proper time, , with τ parametrising proper time along the curve and making manifest the presence of the Christoffel symbols.
A principal feature of general relativity is to determine the paths of particles and radiation in gravitational fields. This is accomplished by solving the geodesic equations.
The EFE relate the total matter (energy) distribution to the curvature of spacetime. Their nonlinearity leads to a problem in determining the precise motion of matter in the resultant spacetime. For example, in a system composed of one planet orbiting a star, the motion of the planet is determined by solving the field equations with the energy–momentum tensor the sum of that for the planet and the star. The gravitational field of the planet affects the total spacetime geometry and hence the motion of objects. It is therefore reasonable to suppose that the field equations can be used to derive the geodesic equations.
When the energy–momentum tensor for a system is that of dust, it may be shown by using the local conservation law for the energy–momentum tensor that the geodesic equations are satisfied exactly.
Lagrangian formulation
The issue of deriving the equations of motion or the field equations in any physical theory is considered by many researchers to be appealing. A fairly universal way of performing these derivations is by using the techniques of variational calculus, the main objects used in this being Lagrangians.
Many consider this approach to be an elegant way of constructing a theory, others as merely a formal way of expressing a theory (usually, the Lagrangian construction is performed after the theory has been developed).
Mathematical techniques for analysing spacetimes
Having outlined the basic mathematical structures used in formulating the theory, some important mathematical techniques that are employed in investigating spacetimes will now be discussed.
Frame fields
A frame field is an orthonormal set of 4 vector fields (1 timelike, 3 spacelike) defined on a spacetime. Each frame field can be thought of as representing an observer in the spacetime moving along the integral curves of the timelike vector field. Every tensor quantity can be expressed in terms of a frame field, in particular, the metric tensor takes on a particularly convenient form. When allied with coframe fields, frame fields provide a powerful tool for analysing spacetimes and physically interpreting the mathematical results.
Symmetry vector fields
Some modern techniques in analysing spacetimes rely heavily on using spacetime symmetries, which are infinitesimally generated by vector fields (usually defined locally) on a spacetime that preserve some feature of the spacetime. The most common type of such symmetry vector fields include Killing vector fields (which preserve the metric structure) and their generalisations called generalised Killing vector fields. Symmetry vector fields find extensive application in the study of exact solutions in general relativity and the set of all such vector fields usually forms a finite-dimensional Lie algebra.
The Cauchy problem
The Cauchy problem (sometimes called the initial value problem) is the attempt at finding a solution to a differential equation given initial conditions. In the context of general relativity, it means the problem of finding solutions to Einstein's field equations a system of hyperbolic partial differential equations given some initial data on a hypersurface. Studying the Cauchy problem allows one to formulate the concept of causality in general relativity, as well as 'parametrising' solutions of the field equations. Ideally, one desires global solutions, but usually local solutions are the best that can be hoped for. Typically, solving this initial value problem requires selection of particular coordinate conditions.
Spinor formalism
Spinors find several important applications in relativity. Their use as a method of analysing spacetimes using tetrads, in particular, in the Newman–Penrose formalism is important.
Another appealing feature of spinors in general relativity is the condensed way in which some tensor equations may be written using the spinor formalism. For example, in classifying the Weyl tensor, determining the various Petrov types becomes much easier when compared with the tensorial counterpart.
Regge calculus
Regge calculus is a formalism which chops up a Lorentzian manifold into discrete 'chunks' (four-dimensional simplicial blocks) and the block edge lengths are taken as the basic variables. A discrete version of the Einstein–Hilbert action is obtained by considering so-called deficit angles of these blocks, a zero deficit angle corresponding to no curvature. This novel idea finds application in approximation methods in numerical relativity and quantum gravity, the latter using a generalisation of Regge calculus.
Singularity theorems
In general relativity, it was noted that, under fairly generic conditions, gravitational collapse will inevitably result in a so-called singularity. A singularity is a point where the solutions to the equations become infinite, indicating that the theory has been probed at inappropriate ranges.
Numerical relativity
Numerical relativity is the sub-field of general relativity which seeks to solve Einstein's equations through the use of numerical methods. Finite difference, finite element and pseudo-spectral methods are used to approximate the solution to the partial differential equations which arise. Novel techniques developed by numerical relativity include the excision method and the puncture method for dealing with the singularities arising in black hole spacetimes. Common research topics include black holes and neutron stars.
Perturbation methods
The nonlinearity of the Einstein field equations often leads one to consider approximation methods in solving them. For example, an important approach is to linearise the field equations. Techniques from perturbation theory find ample application in such areas.
| Physical sciences | Theory of relativity | Physics |
1790574 | https://en.wikipedia.org/wiki/Ecological%20restoration | Ecological restoration | Ecological restoration, or ecosystem restoration, is the process of assisting the recovery of an ecosystem that has been degraded, damaged, destroyed or transformed. It is distinct from conservation in that it attempts to retroactively repair already damaged ecosystems rather than take preventative measures. Ecological restoration can reverse biodiversity loss, combat climate change, support the provision of ecosystem services and support local economies. The United Nations has named 2021-2030 the Decade on Ecosystem Restoration.
Habitat restoration involves the deliberate rehabilitation of a specific area to reestablish a functional ecosystem. This may differ from historical baselines (the ecosystem's original condition at a particular point in time). To achieve successful habitat restoration, it is essential to understand the life cycles and interactions of species, as well as the essential elements such as food, water, nutrients, space, and shelter needed to support species populations.
Scientists estimate that the current species extinction rate, or the rate of the Holocene extinction, is 1,000 to 10,000 times higher than the normal, background rate. Habitat loss is a leading cause of species extinctions and ecosystem service decline. Two methods have been identified to slow the rate of species extinction and ecosystem service decline: conservation of quality habitat and restoration of degraded habitat. The number and size of ecological restoration projects have increased exponentially in recent years.
Restoration goals reflect political choices, and differ by place and culture. On a global level, the concept of nature-positive has emerged as a societal goal to achieve full nature recovery by 2050, including through restoration of degraded ecosystems to reverse biodiversity loss.
Definition
The Society for Ecological Restoration defines restoration as "the process of assisting the recovery of an ecosystem that has been degraded, damaged, or destroyed." Restoration ecology is the academic study of the science of restoration, whereas ecological restoration is the implementation by practitioners. Ecological restoration includes a wide diversity of methods including erosion control, reforestation, removal of non-native species and weeds, revegetation of disturbed areas, daylighting streams, the reintroduction of native species, habitat and range improvement for targeted species and establishing wildlife corridors. Many scholars and practitioners argue that ecological restoration must include local communities and stakeholders: they call this process the "social-ecological restoration".
The goal of ecosystem restoration depends on the specific context of each location. Traditionally, the aim has been to return ecosystems to a past state (historic baseline), based on the idea that past conditions represent a 'pristine' or ideal functioning state. However, this approach is now questioned because human-driven environmental changes, including climate change, continuously alter ecosystems, resulting in a shifting baseline. Today, it's widely recognized that there may be several possible targets for restoration, based on a range of factors. Targets are set based on factors such as the level of ecosystem degradation, how much ecosystem functionality can realistically be restored, local community views, and the costs of restoration efforts.
Rationale
There are many reasons to restore ecosystems. Some include:
Restoring natural capital such as drinkable water or wildlife populations
Restoring environmental degradation
Helping human communities and the ecosystems upon which they depend adapt to the impacts of climate change (through ecosystem-based adaptation)
Mitigating climate change (e.g. through carbon sequestration)
Helping threatened or endangered species
Aesthetic reasons
Moral reasons: human intervention has unnaturally destroyed many habitats, and there exists an innate obligation to restore these destroyed habitats
Regulated use/harvest, particularly for subsistence
Cultural importance to indigenous people
The environmental health of nearby populations
There are considerable differences of opinion on how to set restoration goals and define their success. As Laura J. Martin writes, "Restoration targets are moral and political matters as well as logistical and scientific ones." Some restorationists urge active restoration (e.g. killing invasive animals) and others believe that protected areas should have the bare minimum of human interference, such as rewilding.
Skeptics doubt that the benefits justify the economic investment or point to failed restoration projects and question the feasibility of restoration altogether. It can be difficult to set restoration goals because, as Anthony Bradshaw writes, "ecosystems are not static, but in a state of dynamic equilibrium." Some scientists argue that, though an ecosystem may not be returned to its original state, the functions of a "novel ecosystem" are still valuable.
Ecosystem restoration can mitigate climate change through activities such as afforestation. However, afforestation can have negative impacts on biodiversity especially when considering tree-planting initiatives in tropical savannas. The impacts of afforestation on water supply and quality are also debated and vary by region, climate and age of afforestation projects. Forestry-based carbon offsetting is controversial and sometimes critiqued as carbon colonialism. Another driver of restoration projects in the United States is the legal framework of the Clean Water Act, which often requires mitigation for damage inflicted on aquatic systems by development or other activities.
Theoretical foundations
Ecological restoration draws on a wide range of ecological concepts.
Disturbance
Disturbance is a change in environmental conditions that disrupt the functioning of an ecosystem. Disturbance can occur at a variety of spatial and temporal scales, and is a natural component of many communities. For example, many forest and grassland restorations implement fire as a natural disturbance regime. However the severity and scope of anthropogenic impact has grown in the last few centuries. Differentiating between human-caused and naturally occurring disturbances is important if we are to understand how to restore natural processes and minimize anthropogenic impacts on the ecosystems.
Succession
Ecological succession is the process by which a community changes over time, especially following a disturbance. In many instances, an ecosystem will change from a simple level of organization with a few dominant pioneer species to an increasingly complex community with many interdependent species. Restoration often consists of initiating, assisting, or accelerating ecological successional processes, depending on the severity of the disturbance. Following mild to moderate natural and anthropogenic disturbances, restoration in these systems involves hastening natural successional trajectories through careful management. However, in a system that has experienced a more severe disturbance (such as in urban ecosystems), restoration may require intensive efforts to recreate environmental conditions that favor natural successional processes.
Fragmentation
Habitat fragmentation describes spatial discontinuities in a biological system, where ecosystems are broken up into smaller parts through land-use changes (e.g. agriculture) and natural disturbance. This both reduces the size of the population and increases the degree of isolation. These smaller and isolated populations tend to be more vulnerable to extinction. Fragmenting ecosystems decreases the quality of the habitat. The edge of a fragment has a different range of environmental conditions and therefore supports different species than the interior. Restorative projects can increase the effective size of a population by adding suitable habitat and decrease isolation by creating habitat corridors that link isolated fragments. Reversing the effects of fragmentation is an important component of restoration ecology. The composition of the surrounding landscape can also influence the effectiveness of restoration projects. For example, a restoration site that is closer to remaining vegetation will be more likely to be naturally regenerated through seed disperal than a site that is further away.
Ecosystem function
Ecosystem function describes the most basic and essential foundational processes of any natural systems, including nutrient cycles and energy fluxes. An understanding of the complexity of these ecosystem functions is necessary to address any ecological processes that may be degraded. Ecosystem functions are emergent properties of the system as a whole, thus monitoring and management are crucial for the long-term stability of ecosystems. A completely self-perpetuating and fully functional ecosystem is the ultimate goal of restorative efforts. We must understand what ecosystem properties influence others to restore desired functions and reach this goal.
Community assembly
Community assembly "is a framework that can unify virtually all of (community) ecology under a single conceptual umbrella". Community assembly theory attempts to explain the existence of environmentally similar sites with differing assemblages of species. It assumes that species have similar niche requirements, so that community formation is a product of random fluctuations from a common species pool. Essentially, if all species are fairly ecologically equivalent, then random variation in colonization, and migration and extinction rates between species, drive differences in species composition between sites with comparable environmental conditions.
Population genetics
Genetic diversity has shown to be as important as species diversity for restoring ecosystem processes. Hence ecological restorations are increasingly factoring genetic processes into management practices. Population genetic processes that are important to consider in restored populations include founder effects, inbreeding depression, outbreeding depression, genetic drift, maladaption and gene flow. Such processes can predict whether or not a species successfully establishes at a restoration site.
Applications
Leaf litter accumulation
Leaf litter accumulation plays an important role in the restoration process. Higher quantities of leaf litter hold higher humidity levels, a key factor for the establishment of plants. The process of accumulation depends on factors like wind and species composition of the forest. The leaf litter found in primary forests is more abundant, deeper, and holds more humidity than in secondary forests. These technical considerations are important to take into account when planning a restoration project.
Soil heterogeneity effects on community heterogeneity
Spatial heterogeneity of resources can influence plant community composition, diversity, and assembly trajectory. Baer et al. (2005) manipulated soil resource heterogeneity in a tallgrass prairie restoration project. They found increasing resource heterogeneity, which on its own was insufficient to ensure species diversity in situations where one species may dominate across the range of resource levels. Their findings were consistent with the theory regarding the role of ecological filters on community assembly. The establishment of a single species, best adapted to the physical and biological conditions can play an inordinately important role in determining the community structure.
Invasion and restoration
Restoration is used as a tool for reducing the spread of invasive plant species many ways. The first method views restoration primarily as a means to reduce the presence of invasive species and limit their spread. As this approach emphasizes the control of invaders, the restoration techniques can differ from typical restoration projects. The goal of such projects is not necessarily to restore an entire ecosystem or habitat. These projects frequently use lower diversity mixes of aggressive native species seeded at high density. They are not always actively managed following seeding. The target areas for this type of restoration are those which are heavily dominated by invasive species. The goals are to first remove the species and then in so doing, reduce the number of invasive seeds being spread to surrounding areas. An example of this is through the use of biological control agents (such as herbivorous insects) which suppress invasive weed species while restoration practitioners concurrently seed in native plant species that take advantage of the freed resources. These approaches have been shown to be effective in reducing weeds, although it is not always a sustainable solution long term without additional weed control, such as mowing, or re-seeding.
Restoration projects are also used as a way to better understand what makes an ecological community resistant to invasion. As restoration projects have a broad range of implementation strategies and methods used to control invasive species, they can be used by ecologists to test theories about invasion. Restoration projects have been used to understand how the diversity of the species introduced in the restoration affects invasion. We know that generally higher diversity prairies have lower levels of invasion. The incorporation of functional ecology has shown that more functionally diverse restorations have lower levels of invasion. Furthermore, studies have shown that using native species functionally similar to invasive species are better able to compete with invasive species. Restoration ecologists have also used a variety of strategies employed at different restoration sites to better understand the most successful management techniques to control invasion. To develop restoration ecology into a full science and to improve its practice requires generalizations about the processes governing the development of restored communities. While new experiments can be designed, one way forward is to use data from existing restoration studies to relate plant species performance to their ecological trait.
Successional trajectories
Progress along a desired successional pathway may be difficult if multiple stable states exist. Looking over 40 years of wetland restoration data, Klötzli and Gootjans (2001) argue that unexpected and undesired vegetation assemblies "may indicate that environmental conditions are not suitable for target communities". Succession may move in unpredicted directions, but constricting environmental conditions within a narrow range may rein in the possible successional trajectories and increase the likelihood of the desired outcome.
Sourcing land for restoration
A study quantified climate change mitigation potentials of 'high-income' nations shifting diets – away from meat-consumption – and restoration of the spared land. They find that the hypothetical dietary change "could reduce annual agricultural production emissions of high-income nations' diets by 61% while sequestering as much as 98.3 (55.6–143.7) Gt equivalent, equal to approximately 14 years of current global agricultural emissions until natural vegetation matures", outcomes they call "double climate dividend".
Sourcing material for restoration
For most restoration projects it is generally recommended to source material from local populations, to increase the chance of restoration success and minimize the effects of maladaptation. However the definition of local can vary based on species, habitat and region. US Forest Service recently developed provisional seed zones based on a combination of minimum winter temperature zones, aridity, and the Level III ecoregions. Rather than putting strict distance recommendations, other guidelines recommend sourcing seeds to match similar environmental conditions that the species is exposed to, either now, or under projected climate change. For example, sourcing for Castilleja levisecta found that farther source populations that matched similar environmental variables were better suited for the restoration project than closer source populations. Similarly, a suite of new methods are surveying gene-environment interactions in order to identify the optimum source populations based on genetic adaptation to environmental conditions.
Challenges
Some view ecosystem restoration as impractical, partially because restorations often fall short of their goals. Hilderbrand et al. point out that many times uncertainty (about ecosystem functions, species relationships, and such) is not addressed, and that the time-scales set out for 'complete' restoration are unreasonably short, while other critical markers for full-scale restoration are either ignored or abridged due to feasibility concerns. In other instances an ecosystem may be so degraded that abandonment (allowing a severely degraded ecosystem to recover on its own) may be the wisest option. Local communities sometimes object to restorations that include the introduction of large predators or plants that require disturbance regimes such as regular fires, citing threat to human habitation in the area. High economic costs can also be perceived as a negative impact of the restoration process.
Public opinion is very important in the feasibility of a restoration; if the public believes that the costs of restoration outweigh the benefits they will not support it.
Many failures have occurred in past restoration projects, many times because clear goals were not set out as the aim of the restoration, or an incomplete understanding of the underlying ecological framework lead to insufficient measures. This may be because, as Peter Alpert says, "people may not [always] know how to manage natural systems effectively". Furthermore, many assumptions are made about myths of restoration such as carbon copy, where a restoration plan, which worked in one area, is applied to another with the same results expected, but not realized.
Science–practice gap
One of the struggles for both fields is a divide between restoration ecology and ecological restoration in practice. Many restoration practitioners as well as scientists feel that science is not being adequately incorporated into ecological restoration projects. In a 2009 survey of practitioners and scientists, the "science-practice gap" was listed as the second most commonly cited reason limiting the growth of both science and practice of restoration.
There are a variety of theories about the cause of this gap. However, it has been well established that one of the main issues is that the questions studied by restoration ecologists are frequently not found useful or easily applicable by land managers. For instance, many publications in restoration ecology characterize the scope of a problem in-depth, without providing concrete solutions. Additionally many restoration ecology studies are carried out under controlled conditions and frequently at scales much smaller than actual restorations. Whether or not these patterns hold true in an applied context is often unknown. There is evidence that these small-scale experiments inflate type II error rates and differ from ecological patterns in actual restorations. One approach to addressing this gap has been the development of International Principles & Standards for the Practice of Ecological Restoration by the Society for Ecological restoration (see below) – however, this approach is contended, with scientists active in the field suggesting that this is restrictive, and instead principles and guidelines offer flexibility.
There is further complication in that restoration ecologists who want to collect large-scale data on restoration projects can face enormous hurdles in obtaining the data. Managers vary in how much data they collect, and how many records they keep. Some agencies keep only a handful of physical copies of data that make it difficult for the researcher to access. Many restoration projects are limited by time and money, so data collection and record-keeping are not always feasible. However, this limits the ability of scientists to analyze restoration projects and give recommendations based on empirical data.
Food security and nature degradation
Agriculture is a driver of environmental degradation. However it is vital that ecosystem restoration efforts do not clash with increasing needs for food production. Restoration frameworks aim to assist policy decisions by minimizing trade-offs between ecological restoration and production and evaluating the best use of land to balance carbon storage and food growing. For example, agroforestry is increasing considered as a viable ecosystem restoration strategy, especially in countries with large agriculture footprints.
Restoration as a substitute for steep emission reductions
Climate benefits from nature restoration are "dwarfed by the scale of ongoing fossil fuel emissions". It risks "over-relying on land for mitigation at the expense of phasing out fossil fuels". Despite these issues, nature restoration is receiving increasing attention, with a study concluding that "Land restoration is an important option for tackling climate change but cannot compensate for delays in reducing fossil fuel emissions" as it is "unlikely to be done quickly enough to notably reduce the global peak temperatures expected in the next few decades".
Researchers have found that, in terms of environmental services, it is better to avoid deforestation than to allow for deforestation to subsequently reforest, as the former leads to irreversible effects in terms of biodiversity loss and soil degradation. Furthermore, the probability that legacy carbon will be released from soil is higher in younger boreal forest. Global greenhouse gas emissions caused by damage to tropical rainforests may have been substantially underestimated until around 2019. Additionally, the effects of af- or reforestation will be farther in the future than keeping existing forests intact. It takes much longer − several decades − for the benefits for global warming to manifest to the same carbon sequestration benefits from mature trees in tropical forests and hence from limiting deforestation. Therefore, scientists consider "the protection and recovery of carbon-rich and long-lived ecosystems, especially natural forests" to be "the major climate solution".
Contrasting restoration ecology and conservation biology
Both restoration ecologists and conservation biologists agree that protecting and restoring habitat is important for protecting biodiversity. However, conservation biology is primarily rooted in population biology. Because of that, it is generally organized at the population genetic level and assesses specific species populations (i.e. endangered species). Restoration ecology is organized at the community level, which focuses on broader groups within ecosystems.
In addition, conservation biology often concentrates on vertebrate and invertebrate animals because of their salience and popularity, whereas restoration ecology concentrates on plants. Restoration ecology focuses on plants because restoration projects typically begin by establishing plant communities. Ecological restoration, despite being focused on plants, may also have "umbrella species" for individual ecosystems and restoration projects. For example, the Monarch butterfly is an umbrella species for conserving and restoring milkweed plant habitat, because Monarch butterflies require milkweed plants to reproduce. Finally, restoration ecology has a stronger focus on soils, soil structure, fungi, and microorganisms because soils provide the foundation of functional terrestrial ecosystems.
International Principles & Standards for the Practice of Ecological Restoration
The Society for Ecological Restoration (SER) released the second edition of the International Standards for the Practice of Ecological Restoration on September 27, 2019, in Cape Town, South Africa, at SER's 8th World Conference on Ecological Restoration. The publication provides updated and expanded guidance on the practice of ecological restoration, clarifies the breadth of ecological restoration and allied environmental repair activities, and includes ideas and input from a diverse international group of restoration scientists and practitioners.
The second edition builds on the first edition of the Standards, which was released December 12, 2016, at the Convention on Biological Diversity's 13th Conference of the Parties in Cancun, Mexico. The development of these Standards has been broadly consultative. The first edition was circulated to dozens of practitioners and experts for feedback and review. After release of the first edition, SER held workshops and listening sessions, sought feedback from key international partners and stakeholders, opened a survey to members, affiliates and supporters, and considered and responded to published critiques.
The International Principles and Standards for the Practice of Ecological Restoration:
Present a robust framework to guide restoration projects toward achieving intended goals.
Address restoration challenges including: effective design and implementation, accounting for complex ecosystem dynamics (especially in the context of climate change), and navigating trade-offs associated with land management priorities and decisions.
Highlight the role of ecological restoration in connecting social, community, productivity, and sustainability goals.
Recommend performance measures for restorative activities for industries, communities, and governments to consider.
Enhance the list of practices and actions that guide practitioners in planning, implementation, and monitoring activities, including: appropriate approaches to site assessment and identification of reference ecosystems, different restoration approaches including natural regeneration, and the role of ecological restoration in global restoration initiatives.
Include an expanded glossary of restoration terminology.
Provide a technical appendix on sourcing of seeds and other propagules for restoration.
Implementation by country/region
Indigenous peoples, land managers, stewards, and laypeople have been practicing ecological restoration or ecological management for thousands of years. Restoration ecology emerged as a separate field in ecology in the late twentieth century. The term was coined by John Aber and William Jordan III when they were at the University of Wisconsin–Madison.
European Union
In 2024, the European Union passed a nature restoration law aiming to restore 20% of degraded ecosystems by 2030 and 100% by 2050. The representative of Austria, Leonore Gewessler, voted against the will of its government and can face up to 10 years in prison for doing so.
US
Prior to the emergence of ecology as a scientific discipline, large-scale restoration began with big game restoration in the early 20th century. The first native plant restoration project in the United States was established in 1907 by Eloise Butler in Minneapolis, Minnesota. This was followed by the Vassar College Ecological Laboratory restoration program, founded by Professor Edith Roberts in 1921. The first tallgrass prairie restoration was the 1936 Curtis Prairie at the University of Wisconsin–Madison Arboretum. Civilian Conservation Corps workers replanted nearby prairie species onto a former horse pasture, overseen by university faculty including Aldo Leopold, Theodore Sperry, Henry C. Greene, and John T. Curtis. The UW Arboretum was the center of tallgrass prairie research through the first half of the 20th century and the study of techniques like prescribed burning. It was followed by the 40-hectare Schulenberg Prairie at the Morton Arboretum, initiated in 1962 by Ray Schulenberg and Robert Betz. Betz then worked with The Nature Conservancy to establish the 260-hectare Fermi National Laboratory tallgrass prairie in 1974. Restoration ecology emerged as a distinct sub-discipline of ecology and natural resources management with the dramatic increase in the number of protected natural areas in the 1980s. In 1997 the National Wildlife Federation signed a memorandum of understanding with the Intertribal Bison Cooperative, the first-ever conservation agreement between an environmental organization and an inter-tribal group, to advocate for the restoration of wild bison to tribal lands. Anishinaabek/Neshnabék throughout the Great Lakes region are leading ecological restoration projects that, in the words of Kyle Whyte, "seek to learn from, adapt, and put into practice local human and nonhuman relationships and stories at the convergence of deep Anishinaabe history and the disruptiveness of industrial settler campaigns."
Australia
Australia has been the site of historically significant ecological restoration projects, commencing in the 1930s. These projects were responses to the extensive environmental damage inflicted by colonising settlers, following the forced dispossession of the First Nations communities of Australia. The substantial Traditional Ecological Knowledge of First Nations communities was not utilised in the historical restoration projects.
Many of the first Australian settler restoration projects were initiated by volunteers, often in the form of community groups. Many of these volunteers appreciated and utilised science resources, such as botanical and ecological knowledge. Local and state government agencies participated, and also industry. Australian scientists came to play an increasingly important role. A prominent scientist who took an interest in the reversal of vegetation degradation was botanist and plant ecologist Professor T G Osborn, University of Adelaide, who, in the 1920s, conducted pioneering research into the causes of arid-zone indigenous vegetation degradation. From this time, Australian botanists, plant ecologists and soil erosion researchers have increasingly developed interests in the recovery of ecological functioning on degraded sites.
The earliest known attempt by Australian settlers to restore a degraded natural ecosystem commenced in 1896, at Nairm (as it is known to people of the Kulin nation), or Port Phillip Bay, Melbourne. Local government and community groups replanted degraded areas of the foreshore reserves with the indigenous plant species, coastal teatree (Leptospermum laevigatum). The projects were motivated by utilitarian considerations: to conserve recreation sites, and promote tourism. However, some local residents, including Australian journalist, nature writer and amateur ornithologist, Donald Macdonald, were distressed at the loss of valued biological qualities, and campaigned to fully restore the Teatree ecosystems and conserve them and their indigenous fauna.
The degraded arid-zone regions of Australia were the site of historical ecological restoration projects. Pastoral industry established in the arid-zone regions of South Australia and New South Wales resulted in the substantial degradation of these areas by ca.1900 resulting in severe wind erosion. From approximately 1930, Australian pastoralists implemented revegetation projects aiming to the substantial to full restoration of indigenous flora to degraded, wind eroded areas.
At his arid-zone Koonamore research station in South Australia, established in 1925, Professor T G Osborn studied the loss of indigenous vegetation caused by overstocking and the resultant wind erosion and degradation, concluding that restoration of the indigenous saltbushes (Atriplex spp.), bluebushes (Maireana spp.) and mulga (Acacia aneura) vegetation communities was possible, if a stock exclosure and natural regeneration revegetation technique was applied to degraded paddocks. Most likely influenced by Osborn's research, throughout the 1930s South Australian pastoralists adopted this revegetation technique. For example, at Wirraminna station (or property, ranch), following fencing to exclude stock, severe soil-drifts were fully revegetated and stabilised through natural regeneration of the indigenous vegetation. It was also found that furrowing (or ploughing) of eroded areas resulted in the natural regeneration of indigenous vegetation. So successful were these programs that the South Australian government adopted them as approved state soil conservation policies in 1936. Legislation introduced in 1939 codified these policies.
In 1936 mining assayer Albert Morris and his restoration colleagues initiated the Broken Hill regeneration area project. This project involved the natural regeneration of indigenous flora on a severely wind eroded site of hundreds of hectares, located in arid western New South Wales. Local and state governments, and the Broken Hill mining industry, supported and funded the project. In fact, as the regeneration area project was so well adapted to the harsh arid-zone conditions, the New South Wales state government adopted it as a model for the proposed restoration of the twenty million hectares of the arid western portion of the state that had been reduced to a severely eroded condition. Legislation to this effect was passed in 1949.
Another significant early Australian settler ecological restoration project occurred on the north coast of New South Wales. From approximately 1840 settlers forcibly occupied the coastal hinterlands, dispossessed First Nations communities, destroyed extensive areas of biologically diverse rainforest and converted the land to farms. Only small patches of rainforest survived. In 1935 dairy farmer Ambrose Crawford began restoring a degraded four acre (1.7 hectare) patch of local rainforest, or "Big Scrub" (Lowland Tropical Rainforest), as it was referred to, at Lumley Park reserve, Alstonville. His main restoration techniques were clearing weeds that were smothering the rainforest plants and planting of suitable indigenous rainforest species. Crawford utilised professional government botanists as advisors, and received support from his local government council. The restored rainforest reserve still exists today.
United Kingdom
Natural Capital Committee's recommendation for a 25-year plan
The UK Natural Capital Committee (NCC) made a recommendation in its second State of Natural Capital report published in March 2014 that in order to meet the Government's goal of being the first generation to leave the environment in a better state than it was inherited, a long-term 25-year plan was needed to maintain and improve England's natural capital.
The Secretary of State for the UK's Department for Environment, Food and Rural Affairs, Owen Paterson, described his ambition for the natural environment and how the work of the Committee fits into this at an NCC event in November 2012: "I do not, however, just want to maintain our natural assets; I want to improve them. I want us to derive the greatest possible benefit from them, while ensuring that they are available for generations to come. This is what the NCC's innovative work is geared towards".
Traditional ecological knowledge
Traditional ecological knowledge (TEK) from Indigenous Peoples demonstrates how restoration ecology is a historical field, lived out by humans for thousands of years. Indigenous people have acquired ecological knowledge through observation, experience, and management of the natural resources and the environment around them. In the past, they managed their environment and changed the structure of the vegetation to not only meet their basic needs (food, water, shelter, medicines) but also to improve desired characteristics and even increasing the populations and biodiversity. In that way, they achieved a close relationship with the environment and learned lessons that indigenous people keep in their culture.
This means there is much that could be learned from local people indigenous to the ecosystem being restored because of the deep connection and biocultural and linguistic diversity of place. The use of natural resources by indigenous people considers many cultural, social, and environmental aspects, since they have always had an intimate connection with the animals and plants around them over centuries since they obtained their livelihood from the environment around them.
Restoration ecologists must consider that TEK is place dependent due to intimate connection and thus when engaging Indigenous Peoples to include knowledge for restoration purposes, respect and care must be taken to avoid appropriation of the TEK. Successful ecological restoration which includes Indigenous Peoples must be led by Indigenous Peoples to ensure non-indigenous people acknowledge the unequal relationship of power.
For example, the California Indians have a rigid and complex harvesting, management and production practice, largely typical horticultural techniques and concentrated forest burning. The California Indians had a rich knowledge of ecology and natural techniques to understand burn patterns, plant material, cultivation, pruning, digging; what was edible vs. what was not. This knowledge extends into wildlife management – how abundant, where the distribution was, and how diverse the large mammal population was. While the United States has counteracted the degradation, fragmentation and loss of habitat through land set aside from all human influence, indigenous practices could inform ecosystem restoration and wildlife management.
Related journals
Restoration Ecology, journal of the Society for Ecological Restoration (SER)
Ecological Management & Restoration, published by the Ecological Society of Australia (ESA)
Ecological Restoration, published by the University of Wisconsin Press
| Biology and health sciences | Ecology | Biology |
20471890 | https://en.wikipedia.org/wiki/Per-%20and%20polyfluoroalkyl%20substances | Per- and polyfluoroalkyl substances | Per- and polyfluoroalkyl substances (PFAS or PFASs) are a group of synthetic organofluorine chemical compounds that have multiple fluorine atoms attached to an alkyl chain; there are 7 million such chemicals according to PubChem. PFAS came into use after the invention of Teflon in 1938 to make fluoropolymer coatings and products that resist heat, oil, stains, grease, and water. They are now used in products including waterproof fabric such as Nylon, yoga pants, carpets, shampoo, feminine hygiene products, mobile phone screens, wall paint, furniture, adhesives, food packaging, heat-resistant non-stick cooking surfaces such as Teflon, firefighting foam, and the insulation of electrical wire. PFAS are also used by the cosmetic industry in most cosmetics and personal care products, including lipstick, eye liner, mascara, foundation, concealer, lip balm, blush, and nail polish.
Many PFAS such as PFOS and PFOA pose health and environmental concerns because they are persistent organic pollutants; they were branded as "forever chemicals" in an article in The Washington Post in 2018. Some have half-lives of over eight years due to a carbon-fluorine bond, one of the strongest in organic chemistry. They move through soils and bioaccumulate in fish and wildlife, which are then eaten by humans. Residues are now commonly found in rain, drinking water, and wastewater. Since PFAS compounds are highly mobile, they are readily absorbed through human skin and through tear ducts, and such products on lips are often unwittingly ingested. Due to the large number of PFAS, it is challenging to study and assess the potential human health and environmental risks; more research is necessary and is ongoing.
Exposure to PFAS, some of which have been classified as carcinogenic and/or as endocrine disruptors, has been linked to cancers such as kidney, prostate and testicular cancer, ulcerative colitis, thyroid disease, suboptimal antibody response / decreased immunity, decreased fertility, hypertensive disorders in pregnancy, reduced infant and fetal growth and developmental issues in children, obesity, dyslipidemia (abnormally high cholesterol), and higher rates of hormone interference.
The use of PFAS has been regulated internationally by the Stockholm Convention on Persistent Organic Pollutants since 2009, with some jurisdictions, such as China and the European Union, planning further reductions and phase-outs. However, major producers and users such as the United States, Israel, and Malaysia have not ratified the agreement and the chemical industry has lobbied governments to reduce regulations or has moved production to countries such as Thailand, where there is less regulation. In the United States, the Republican Party has filibustered bills regulating the chemicals. Cover-ups and the suppression of studies in 2018 by the Trump administration led to bipartisan outrage.
The market for PFAS was estimated to be $28 billion in 2023 and the majority are produced by 12 companies: 3M, AGC Inc., Archroma, Arkema, BASF, Bayer, Chemours, Daikin, Honeywell, Merck Group, Shandong Dongyue Chemical, and Solvay. Sales of PFAS, which cost approximately $20 per kilogram, generate a total industry profit of $4 billion per year on 16% profit margins. Due to health concerns, several companies have ended or plan to end the sale of PFAS or products that contain them; these include W. L. Gore & Associates (the maker of Gore-Tex), H&M, Patagonia, REI, and 3M. PFAS producers have paid billions of dollars to settle litigation claims, the largest being a $10.3 billion settlement paid by 3M for water contamination in 2023. Studies have shown that companies have known of the health dangers since the 1970s – DuPont and 3M were aware that PFAS was "highly toxic when inhaled and moderately toxic when ingested". External costs, including those associated with remediation of PFAS from soil and water contamination, treatment of related diseases, and monitoring of PFAS pollution, may be as high as US$17.5 trillion annually, according to ChemSec. The Nordic Council of Ministers estimated health costs to be at least €52–84 billion in the European Economic Area. In the United States, PFAS-attributable disease costs are estimated to be US$6–62 billion.
In January 2025, reports stated that the cost of cleaning up toxic PFAS pollution in the UK and Europe could exceed £1.6 trillion over the next 20 years, averaging £84 billion annually.
Definition
Per- and polyfluoroalkyl substances are a group of synthetic organofluorine chemical compounds that have multiple fluorine atoms attached to an alkyl chain. Different organizations use different definitions for PFAS, leading to estimates of between 8,000 and 7 million chemicals within the group. The EPA toxicity database, DSSTox, lists 14,735 unique PFAS chemical compounds.
An early definition required that they contain at least one perfluoroalkyl moiety, . Beginning in 2021, the OECD expanded its terminology, stating that "PFAS are defined as fluorinated substances that contain at least one fully fluorinated methyl or methylene carbon atom (without any H/Cl/Br/I atom attached to it), i.e., with a few noted exceptions, any chemical with at least a perfluorinated methyl group () or a perfluorinated methylene group () is a PFAS." This definition notably includes Carbon tetrafluoride.
The United States Environmental Protection Agency defines PFAS in the Drinking Water Contaminant Candidate List 5 as substances that contain "at least one of the following three structures: , where both the and moieties are saturated carbons, and none of the R groups can be hydrogen; , where both the moieties are saturated carbons, and none of the R groups can be hydrogen; or , where all the carbons are saturated, and none of the R groups can be hydrogen.
A summary table of some PFAS definitions is provided in Hammel et al (2022).
Fluorosurfactants
Fluorinated surfactants or fluorosurfactants are a subgroup of PFAS characterized by a hydrophobic fluorinated "tail" and a hydrophilic "head" that behave as surfactants. These are more effective at reducing the surface tension of water than comparable hydrocarbon surfactants. They include the perfluorosulfonic acids, such as perfluorooctanesulfonic acid (PFOS), and the perfluorocarboxylic acids like perfluorooctanoic acid (PFOA).
As with other surfactants, fluorosurfactants tend to concentrate at the phase interfaces. Fluorocarbons are both lipophobic and hydrophobic, repelling both oil and water. Their lipophobicity results from the relative lack of London dispersion forces compared to hydrocarbons, a consequence of fluorine's large electronegativity and small bond length, which reduce the polarizability of the surfactants' fluorinated molecular surface. Fluorosurfactants are more stable than hydrocarbon surfactants due to the stability of the carbon–fluorine bond. Perfluorinated surfactants persist in the environment for the same reason.
Fluorosurfactants such as PFOS, PFOA, and perfluorononanoic acid (PFNA) have caught the attention of regulatory agencies because of their persistence, toxicity, and widespread occurrence in the blood of general populations.
PFASs are used in emulsion polymerization to produce fluoropolymers, used in stain repellents, polishes, paints, and coatings.
Health and environmental effects
PFASs were originally considered to be chemically inert. Early occupational studies revealed elevated levels of fluorochemicals, including perfluorooctanesulfonic acid (PFOS) and perfluorooctanoic acid (PFOA, C8), in the blood of exposed industrial workers, but cited no ill health effects. These results were consistent with the measured serum concentrations of PFOS and PFOA in 3M plant workers ranging from 0.04 to 10.06 ppm and 0.01 to 12.70 ppm, respectively, well below toxic and carcinogenic levels cited in animal studies. Given, however, the serum elimination half-life of four to five years and widespread environmental contamination, molecules have been shown to accumulate in humans sufficiently to cause adverse health outcomes.
Prevalence in rain, soil, water bodies, and air
In 2022, levels of at least four perfluoroalkyl acids (PFAAs) in rain water worldwide greatly exceeded the EPA's lifetime drinking water health advisories as well as comparable Danish, Dutch, and European Union safety standards, leading to the conclusion that "the global spread of these four PFAAs in the atmosphere has led to the planetary boundary for chemical pollution being exceeded".
It had been thought that PFAAs would eventually end up in the oceans, where they would be diluted over decades, but a field study published in 2021 by researchers at Stockholm University found that they are often transferred from water to air when waves reach land, are a significant source of air pollution, and eventually get into rain. The researchers concluded that pollution may impact large areas.
In 2024, a worldwide study of 45,000 groundwater samples found that 31% of samples contained levels of PFAS that were harmful to human health; these samples were taken from areas not near any obvious source of contamination.
Soil is also contaminated and the chemicals have been found in remote areas such as Antarctica. Soil contamination can result in higher levels of PFAs found in foods such as white rice, coffee, and animals reared on contaminated ground.
Adverse health outcomes
From 2005 to 2013, three epidemiologists known as the C8 Science Panel conducted health studies in the Mid-Ohio Valley as part of a contingency to a class action lawsuit brought by communities in the Ohio River Valley against DuPont in response to landfill and wastewater dumping of PFAS-laden material from DuPont's West Virginia Washington Works plant. The panel measured PFOA (also known as C8) serum concentrations in 69,000 individuals from around DuPont's Washington Works Plant and found a mean concentration of 83 ng/mL, compared to 4 ng/mL in a standard population of Americans. This panel reported probable links between elevated PFOA blood concentration and hypercholesterolemia, ulcerative colitis, thyroid disease, testicular cancer, kidney cancer as well as pregnancy-induced hypertension and preeclampsia.
The severity of PFAS-associated health effects can vary based on the length of exposure, level of exposure, and health status.
Pregnancy and lactation issues
Exposure to PFAS is a risk factor for various hypertensive disorders in pregnancy, including preeclampsia and high blood pressure. It is not clear whether PFAS exposure is associated with wider cardiovascular disorders during pregnancy. Human breast milk can harbor PFASs, which can be transferred from mother to infant via breastfeeding.
Use of various personal care products, such as nail care products, fragrances, makeup, hair dyes and hair sprays, by pregnant women and lactating mothers has been shown to be associated with significantly higher levels of PFAS in the blood and breastmilk of the mothers. For example, PFOS levels of women who dyed their hair at least twice during pregnancy were more than a third higher than those who did not. PFOS is one of the most common and most dangerous of the PFAS compounds.
Fertility issues
Endocrine disruptors, including PFASs, are linked with the male infertility crisis.
A report in 2023 by the Icahn School of Medicine at Mount Sinai linked high exposure to PFAS with a 40% decrease in the ability for a woman to have a successful pregnancy as well as hormone disruption and delayed puberty onset.
Human developmental issues
Fetuses and children are especially vulnerable to the harms of PFAS chemicals because they have been shown to be linked to major adverse health conditions, including abnormally small birth weight syndrome in newborns, preterm birth, shorter lactation periods, breastmilk of diminished nutritional content, one or more neurodevelopmental disorders, and decreased response to childhood vaccines.
Liver issues
A meta-analysis for associations between PFASs and human clinical biomarkers for liver injury, analyzing PFAS effects on liver biomarkers and histological data from rodent experimental studies, concluded that evidence exists that PFOA, perfluorohexanesulfonic acid (PFHxS), and perfluorononanoic acid (PFNA) caused hepatotoxicity in humans.
Cancer
PFOA is classified as carcinogenic to humans (Group 1) by the International Agency for Research on Cancer (IARC) based on "sufficient" evidence for cancer in animals and "strong" mechanistic evidence in exposed humans. IARC also classified PFOS as possibly carcinogenic to humans (Group 2b) based on "strong" mechanistic evidence. There is a lack of high-quality epidemiological data on the associations between many specific PFAS chemicals and specific cancer types, and research is ongoing.
Hypercholesterolemia
A response is observed in humans where elevated PFOS levels were significantly associated with elevated total cholesterol and LDL cholesterol, highlighting significantly reduced PPAR expression and alluding to PPAR independent pathways predominating over lipid metabolism in humans compared to rodents.
Ulcerative colitis
PFOA and PFOS have been shown to significantly alter immune and inflammatory responses in human and animal species. In particular, IgA, IgE (in females only) and C-reactive protein have been shown to decrease whereas antinuclear antibodies increase as PFOA serum concentrations increase. These cytokine variations allude to immune response aberrations resulting in autoimmunity. One proposed mechanism is a shift towards anti-inflammatory M2 macrophages and/or (TH2) response in intestinal epithelial tissue which allows sulfate-reducing bacteria to flourish. Elevated levels of hydrogen sulfide result, which reduce beta-oxidation and nutrient production, leading to a breakdown of the colonic epithelial barrier.
Thyroid disease
Hypothyroidism is the most common thyroid abnormality associated with PFAS exposure. PFASs have been shown to decrease thyroid peroxidase, resulting in decreased production and activation of thyroid hormones in vivo. Other proposed mechanisms include alterations in thyroid hormone signaling, metabolism and excretion as well as function of nuclear hormone receptor.
Bioaccumulation and biomagnification
In marine species of the food web
Bioaccumulation controls internal concentrations of pollutants, including PFAS, in individual organisms. When bioaccumulation is looked at in the perspective of the entire food web, it is called biomagnification, which is important to track because lower concentrations of pollutants in environmental matrices such as seawater or sediments, can very quickly grow to harmful concentrations in organisms at higher trophic levels, including humans. Notably, concentrations in biota can even be greater than 5000 times those present in water for PFOS and C10–C14 PFCAs. PFAS can enter an organism by ingestion of sediment, through the water, or directly via their diet. It accumulates namely in areas with high protein content, in the blood and liver, but it is also found to a lesser extent in tissues.
Biomagnification can be described using the estimation of the trophic magnification factor (TMF), which describes the relationship between the contamination levels in a species and their trophic level in the food web. TMFs are determined by graphing the log-transformed concentrations of PFAS against the assigned trophic level and taking the antilog of the regression slope (10slope).
In a study done on a macrotidal estuary in Gironde, SW France, TMFs exceeded one for nearly all 19 PFAS compounds considered in the study and were particularly high for PFOA and PFNA (6.0 and 3.1 respectively). A TMF greater than one signifies that the concentration of a chemical in organisms increases at successive trophic levels, thereby demonstrating biomagnification.
PFOS, a long-chain sulfonic acid, was found at the highest concentrations relative to other PFASs measured in fish and birds in northern seas such as the Barents Sea and the Canadian Arctic.
A study published in 2023 analyzing 500 composite samples of fish fillets collected across the United States from 2013 to 2015 under the EPA's monitoring programs showed freshwater fish ubiquitously contain high levels of harmful PFAS, with a single serving typically significantly increasing the blood PFOS level.
Bioaccumulation and biomagnification of PFASs in marine species throughout the food web, particularly frequently consumed fish and shellfish, can have important impacts on human populations. PFASs have been frequently documented in both fish and shellfish that are commonly consumed by human populations, which poses health risks to humans and studies on the bioaccumulation in certain species are important to determine daily tolerable limits for human consumption, and where those limits may be exceeded causing potential health risks. This has particular implications for populations that consume larger numbers of wild fish and shellfish species. PFAS contamination has also resulted in disruptions to the food supply, such as closures and limits on fishing.
Fluorosurfactants with shorter carbon chains may be less prone to accumulating in mammals; there is still some concern that they may be harmful to both humans and the environment.
Suppression of information on health effects
Since the 1970s, DuPont and 3M were aware that PFAS was "highly toxic when inhaled and moderately toxic when ingested". Producers used several strategies to influence science and regulation – most notably, suppressing unfavorable research and distorting public discourse.
In 2018, under the Presidency of Donald Trump, White House staff and the United States Environmental Protection Agency pressured the U.S. Agency for Toxic Substances and Disease Registry to suppress a study that showed PFASs to be even more dangerous than previously thought.
Concerns, litigation, and regulations in specific countries and regions
Arctic
In 2024, research at McGill University in Quebec, indicated that PFASs were being brought to the Arctic from polluted southern waters by migrating birds. Although it is much less than compared to the introduction by wind and the oceans, the birds become vectors, transmitting the toxic chemicals. Rainer Lohmann, an oceanographer at the University of Rhode Island, noted that this has a significant localized affect that is devastating for Arctic predators who accumulate toxins in their bodies because the contaminants from the birds often enter the food chain directly since the birds are the prey of many species.
Australia
In 2017, the ABC's current affairs program Four Corners reported that the storage and use of firefighting foams containing perfluorinated surfactants at Australian Defence Force facilities around Australia had contaminated nearby water resources. In 2019, remediation efforts at RAAF Base Tindal and the adjacent town of Katherine were ongoing. In the 2022 Australian federal budget $428million was allocated for works at HMAS Albatross, RAAF Base Amberley, RAAF Base Pearce and RAAF Base Richmond including funding to remediate PFAS contamination.
Canada
Although PFASs are not manufactured in Canada, they may be present in imported goods and products. In 2008, products containing PFOS as well as PFOA were banned in Canada, with exceptions for products used in firefighting, the military, and some forms of ink and photo media.
Health Canada has published drinking water guidelines for maximum concentrations of PFOS and PFOA to protect the health of Canadians, including children, over a lifetime's exposure to these substances. The maximum allowable concentration for PFOS under the guidelines is 0.0002 milligrams per liter. The maximum allowable concentration for PFOA is 0.0006 milligrams per liter. In August 2024, Health Canada established an objective of 30 ng/L for the sum of the concentration of 25 PFASs detected in drinking water.
New Zealand
The New Zealand Environmental Protection Agency (EPA) has banned the use of per- and polyfluoroalkyl substances (PFAS) in cosmetic products starting from 31 December 2026. This will make the country one of the first in the world to take this step on PFAS to protect people and the environment.
United Kingdom
The environmental consequences of PFAS, especially from firefighting activities, have been recognized since the mid-1990s and came to prominence after the Buncefield explosion on 11 December 2005. The Environment Agency has undertaken a series of projects to understand the scale and nature of PFAS in the environment. The Drinking Water Inspectorate requires water companies to report concentrations of 47 PFAS.
European Union
Many PFASs are either not covered by European legislation or are excluded from registration obligations under the EU Registration, Evaluation, Authorisation and Restriction of Chemicals (REACH) chemical regulation. Several PFASs have been detected in drinking water, municipal wastewater, and landfill leachates worldwide.
In 2019, the European Council requested the European Commission to develop an action plan to eliminate all non-essential uses of PFAS due to the growing evidence of adverse effects caused by exposure to these substances; the evidence for the widespread occurrence of PFAS in water, soil, articles, and waste; and the threat it can pose to drinking water. Germany, the Netherlands, Denmark, Norway, and Sweden submitted a so-called restriction proposal based on the REACH regulation to achieve a European ban on the production, use, sale and import of PFAS. The proposal states that a ban is necessary for all use of PFAS, with different periods for different applications when the ban takes effect (immediately after the restriction comes into force, five years afterward, or 12 years afterward), depending on the function and the availability of alternatives. The proposal has not assessed the use of PFAS in medicines, plant protection products, and biocides because specific regulations apply to those substances (Biocidal Products Regulation, Plant Protection Products Regulation, Medicinal Products Regulation) that have an explicit authorization procedure that focuses on risk for health and the environment.
The proposal was submitted on 13 January 2023 and published by the European Chemicals Agency (ECHA) on 7 February. From 22 March to 21 September, citizens, companies, and other organizations commented on the proposal during a public consultation. Based on the information in the restriction proposal and the consultation, two committees from ECHA formulate an opinion on the risk and socio-economic aspects of the proposed restriction. Within a year of publication, the opinions are sent to the European Commission, which makes a final proposal that is submitted to the EU Member States for discussion and decision. Eighteen months after the publication of the restriction decision (which may differ from the original proposal), it will enter into force.
Italy
127,000 residents in the Veneto region are estimated to have been exposed to contamination through tap water, and it is thought to be Europe's biggest PFAS-related environmental disaster. While Italy's National Health Institute (ISS, Istituto Superiore di Sanità) set the threshold limit of PFOA in the bloodstream at 8 nanograms per milliliter (ng/mL), some residents had reached 262 and some industrial employees reach 91,900 ng/mL. In 2021 some data was disclosed by Greenpeace and local citizens after a long legal battle against the Veneto Region and ISS, which for years has denied access to data, despite values known since or even before 2017. The Veneto region has not carried out further monitoring or taken resolutive actions to eliminate pollution and reduce, at least gradually, the contamination of non-potable water. Although in 2020 the European Food Safety Agency (EFSA) has reduced by more than four times the maximum tolerable limit of PFAS that can be taken through the diet, the region has not carried out new assessments or implemented concrete actions to protect the population and the agri-food and livestock sectors. Some limits were added to monitoring the geographical area, which does not include the orange zone and other areas affected by contamination, as well as the insufficiency of analysis on important productions widespread in the areas concerned: eggs (up to 37,100 ng/kg), fish (18,600 ng/kg) spinach and radicchio (only one sampling carried out), kiwis, melons, watermelons, cereals (only one sample was analyzed), soy, wines and apples.
Japan
A study of public water bodies ending in March 2022 showed that the sum of PFOS and PFOA concentrations exceeded 50 ng/L in 81 out of 1,133 test sites and in some cases are present at elevated levels in blood. This has led to pressure to increase regulations.
Sweden
Highly contaminated drinking water has been detected at several locations in Sweden. Such locations include Arvidsjaur, Lulnäset, Uppsala and Visby. In 2013, PFAS were detected at high concentrations in one of the two municipality drinking water treatment plants in the town of Ronneby, in southern Sweden. Concentrations of PFHxS and PFOS were found at 1700 ng/L and 8000 ng/L, respectively. The source of contamination was later found to be a military fire-fighting exercise site in which PFAS containing fire-fighting foam had been used since the mid-1980s.
Additionally, low-level contaminated drinking water has also been shown to be a significant exposure source of PFOA, PFNA, PFHxS and PFOS for Swedish adolescents (ages 10–21). Even though the median concentrations in the municipality drinking water were below one ng/L for each individual PFAS, positive associations were found between adolescent serum PFAS concentrations and PFAS concentrations in drinking water.
United States
An estimated 26,000 U.S. sites are contaminated with PFASs. More than 200 million Americans are estimated to live in places where the PFAS level in tap water, including PFOA and PFOS levels, exceeds the 1 ppt (part per trillion) limit set in 2022 by the EPA.
Based on tap water studies from 716 locations from 2016 and 2021, the U.S. Geological Survey (USGS) found that the PFAS levels exceeded the EPA advisories in approximately 75% of the samples from urban areas and in approximately 25% of the rural area samples.
Certain PFASs are no longer manufactured in the United States as a result of phase-outs including the PFOA Stewardship Program (2010–2015), in which eight major chemical manufacturers agreed to eliminate the use of PFOA and PFOA-related chemicals in their products and emissions from their facilities. However, they are still produced internationally and are imported into the U.S. in consumer goods. Some types of PFAS are voluntarily not included in food packaging.
In 2021, Senators Susan Collins of Maine and Richard Blumenthal of Connecticut proposed the No PFAS in Cosmetics Act in the United States Senate. It was also introduced in the United States House of Representatives by Michigan Representative Debbie Dingell, but the Republican Party, supported by the U.S. chemical industry filibustered the bill.
Military bases
The water in and around at least 126 U.S. military bases has been contaminated by high levels of PFASs because of their use of firefighting foams since the 1970s, according to a study by the U.S. Department of Defense. Of these, 90 bases reported PFAS contamination that had spread to drinking water or groundwater off the base.
In 2022, a report by the Pentagon acknowledged that approximately 175,000 U.S. military personnel at two dozen American military facilities drank water contaminated by PFAS that exceeded the U.S. EPA limit. However, according to the Environmental Working Group, the Pentagon report downplayed the number of people exposed to PFAS, which was probably over 640,000 at 116 military facilities. The EWG found that the Pentagon also omitted from its report some types of diseases that are likely to be caused by PFAS exposure, such as testicular cancer, kidney disease, and fetal abnormalities.
Environmental Protection Agency actions
The United States Environmental Protection Agency has published non-enforceable drinking water health advisories for PFOA and PFOS. In March 2021 EPA announced that it would develop national drinking water standards for PFOA and PFOS. Drinking water utilities are required to monitor PFAS levels and may receive subsidies to do so. There are also regulations regarding wastewater (effluent guidelines) for industries that use PFASs in the manufacturing process as well as biosolids (processed wastewater sludge used as fertilizer).
The EPA issued health advisories for four specific PFASs in June 2022, significantly lowering their safe threshold levels for drinking water. PFOA was reduced from 70 ppt to 0.004 ppt, while PFOS was reduced from 70 ppt to 0.02 ppt. A safe level for the compound GenX was set at 10 ppt, while that for PFBS was set at 2000 ppt. While not enforceable, these health advisories are intended to be acted on by states in setting their own drinking water standards.
In August 2022, the EPA proposed to add PFOA and PFOS to its list of hazardous substances under the Superfund law. EPA issued a final rule in April 2024, which requires that polluters pay for investigations and cleanup of these substances.
In April 2024, the EPA issued a final drinking water rule for PFOA, PFOS, GenX, PFBS, PFNA, and PFHxS. Within three years, public water systems must remove these six PFAS to near-zero levels. States may be awarded grants up to $1 billion in aid to help with the initial testing and treatment of water for this purpose.
Legal actions
In February 2017, DuPont and Chemours (a DuPont spin-off) agreed to pay $671 million to settle lawsuits arising from 3,550 personal injury claims related to the releasing of PFASs from their Parkersburg, West Virginia, plant into the drinking water of several thousand residents. This was after a court-created independent scientific panel—the C8 Science Panel—found a "probable link" between C8 exposure and six illnesses: kidney and testicular cancer, ulcerative colitis, thyroid disease, pregnancy-induced hypertension and high cholesterol.
In October 2018, a class action suit was filed by an Ohio firefighter against several producers of fluorosurfactants, including 3M and DuPont, on behalf of all U.S. residents who may have adverse health effects from exposure to PFASs. The story is told in the film Dark Waters.
In June 2023, 3M reached a US$10.3 billion settlement with several US public water providers to resolve water pollution claims tied to PFAS, while Chemours, DuPont and Corteva settled similar claims for $1.19 billion.
In December 2023, as part of a four-year legal battle, the EPA banned Inhance, a Houston, Texas-based manufacturer that produces an estimated 200 million containers annually with a process that creates PFOA, from using the manufacturing process. In March 2024, the United States Court of Appeals for the Fifth Circuit overturned the ban. While the court did not deny the containers’ health risks, it said that the EPA could not regulate the manufactured containers under Toxic Substances Control Act of 1976, which only addresses "new" chemicals.
State actions
In 2021, Maine became the first U.S. state to ban these compounds in all products by 2030, except for instances deemed "currently unavoidable".
, the states of California, Connecticut, Massachusetts, Michigan, Minnesota, New Hampshire, New Jersey, New York, Vermont, and Wisconsin had enforceable drinking water standards for between two and six types of PFAS. The six chemicals (termed by the Massachusetts Department of Environmental Protection as PFAS6) are measured either individually or summed as a group depending on the standard; they are:
Perfluorooctanesulfonic acid (PFOS)
Perfluorooctanoic acid (PFOA)
Perfluorohexanesulfonic acid (PFHxS)
Perfluorononanoic acid (PFNA)
Perfluoroheptanoic acid (PFHpA)
Perfluorodecanoic acid (PFDA)
California
In 2021 California banned PFASs for use in food packaging and from infant and children's products and also required PFAS cookware in the state to carry a warning label.
Maine
A program licensed and promoted by the Maine Department of Environmental Protection that provided free municipal wastewater sludge (biosolids) to farmers as fertilizer has resulted in PFAS contamination of local drinking water and farm-grown produce.
Michigan
The Michigan PFAS Action Response Team (MPART) was launched in 2017 and is the first multi-agency action team of its kind in the nation. Agencies representing health, environment, and other branches of state government have joined together to investigate sources and locations of PFAS contamination in the state, take action to protect people's drinking water, and keep the public informed. Groundwater is tested at locations throughout the state by various parties to ensure safety, compliance with regulations, and proactively detect and remedy potential problems. In 2010, the Michigan Department of Environmental Quality (MDEQ) discovered levels of PFASs in groundwater monitoring wells at the former Wurtsmith Air Force Base. In 2024, citizen-led testing near the base in Oscoda discovered high levels of PFAS in foam along the shore of Lake Huron. As additional information became available from other national testing, Michigan expanded its investigations into other locations where PFAS compounds were potentially used. In 2018, the MDEQ's Remediation and Redevelopment Division (RRD) established cleanup criteria for groundwater used as drinking water of 70 ppt of PFOA and PFOS, individually or combined. The RRD staff are responsible for implementing these criteria as part of their ongoing efforts to clean up sites of environmental contamination. The RRD staff are the lead investigators at most of the PFAS sites on the MPART website and also conduct interim response activities, such as coordinating bottled water or filter installations with local health departments at sites under investigation or with known PFAS concerns. Most of the groundwater sampling at PFAS sites under RRD's lead is conducted by contractors familiar with PFAS sampling techniques. The RRD also has a Geologic Services Unit, with staff who install monitoring wells and are also well versed with PFAS sampling techniques. The MDEQ has been conducting environmental clean-up of regulated contaminants for decades. Due to the evolving nature of PFAS regulations as new science becomes available, the RRD is evaluating the need for regular PFAS sampling at Superfund sites and is including an evaluation of PFAS sampling needs as part of a Baseline Environmental Assessment review. Earlier in 2018, the RRD purchased lab equipment that will allow the MDEQ Environmental Lab to conduct analyses of certain PFAS samples. (Currently, most samples are shipped to one of the few labs in the country that conduct PFAS analysis, in California, although private labs in other parts of the country, including Michigan, are starting to offer these services.) As of August 2018, RRD has hired additional staff to work on developing the methodology and conducting PFAS analyses.
In 2020 Michigan Attorney General Dana Nessel filed a lawsuit against 17 companies, including 3M, Chemours, and DuPont, for hiding known health and environmental risks from the state and its residents. Nessel's complaint identifies 37 sites with known contamination. The Michigan Department of Environment, Great Lakes, and Energy introduced some of the strictest drinking water standards in the country for PFAS, setting maximum contaminant levels (MCLs) for PFOA and PFOS to 8 and 16 ppt respectively (down from previous existing groundwater cleanup standards of 70 ppt for both), and introducing MCLs for five other previously unregulated PFAS compounds, limiting PFNA to six ppt, PFHxA to 400,000 ppt, PFHxS to 51 ppt, PFBS to 420 ppt and HFPO-DA to 370 ppt. The change adds 38 additional sites to the state's list of known PFAS contaminated areas, bringing the total number of known sites to 137. About half of these sites are landfills and 13 are former plating facilities.
In 2022 PFOS was found in beef produced at a Michigan farm: the cattle had been fed crops fertilized with contaminated biosolids. State agencies issued a consumption advisory, but did not order a recall, because there currently is no PFOS contamination in beef government standards.
A 2024 study found that "atmospheric deposition could be a significant environmental pathway, particularly for the Great Lakes."
Minnesota
In February 2018, 3M settled a lawsuit for $850 million related to contaminated drinking water in Minnesota.
New Jersey
In 2018 the New Jersey Department of Environmental Protection (NJDEP) published a drinking water standard for PFNA. Public water systems in New Jersey are required to meet an MCL standard of 13 ppt. In 2020 the state set a PFOA standard at 14 ppt and a PFOS standard at 13 ppt.
In 2019 NJDEP filed lawsuits against the owners of two plants that had manufactured PFASs, and two plants that were cited for water pollution from other chemicals. The companies cited are DuPont, Chemours, and 3M. NJDEP also declared five companies to be financially responsible for statewide remediation of the chemicals. Among the companies accused were Arkema and Solvay regarding a West Deptford Facility in Gloucester County, where Arkema manufactured PFASs, but Solvay claims to have never manufactured but only handled PFASs. The companies denied liability and contested the directive. In June 2020, the EPA and New Jersey Department of Environmental Protection published a paper reporting that a unique family of PFAS used by Solvay, chloroperfluoropolyether carboxylates (ClPFPECAs), were contaminating the soils of New Jersey as far from the Solvay facility as 150 km. and the ClPFPECAs were found in water as well.
Later in 2020, the New Jersey state attorney general filed suit in the New Jersey Superior Court against Solvay regarding PFAS contamination of the state's environment. In May 2021, Solvay issued a press release that the company is "discontinuing the use of fluorosurfactants in the U.S.".
New York
In 2016, New York, along with Vermont and New Hampshire, acknowledged PFOA contamination by requesting the EPA to release water quality guidance measures. Contamination has been observed by the New York State Department of Environmental Conservation in Hoosick Falls, Newburgh, Petersburgh, Poestenkill, Mahopac, and Armonk.
After a class action lawsuit, in 2021, the village of Hoosick Falls received a $65.25 million settlement from Saint-Gobain Performance Plastics, Honeywell, 3M, and DuPont due to the disposal of PFAS chemicals into the groundwater of the local water treatment plant.
Washington
Five military installations in Washington State have been identified by the United States Senate Committee on Environment and Public Works as having PFAS contamination. Toward environmental and consumer protections, the Washington State Department of Ecology published a Chemical Action Plan in November 2021, and in June 2022 the governor tasked the Washington State Department of Ecology with phasing out manufacture and import of products containing PFASs. Initial steps taken by the Washington State Department of Health to protect the public from exposure through drinking water have included setting State Action Levels for five PFASs (PFOA, PFOS, PFNA, PFHxS, and PFBS), which were implemented in November 2021.
United Nations
In 2009, PFOS, its salts, and perfluorooctanesulfonyl fluoride, as well as PFOA and PFHxS, were listed as persistent organic pollutants under the Stockholm Convention on Persistent Organic Pollutants due to their ubiquitous, persistent, bioaccumulative, and toxic nature. The convention has been ratified by 186 jurisdictions, but has most notably not been ratified by the United States, Israel, and Malaysia. The long-chain (C9–C21) PFCAs are currently under review for listing.
Occupational exposure
Occupational exposure to PFASs occurs in numerous industries due to the widespread use of the chemicals in products and as an element of industrial process streams. PFASs are used in more than 200 different ways in industries as diverse as electronics and equipment manufacturing, plastic and rubber production, food and textile production, and building and construction. Occupational exposure to PFASs can occur at fluorochemical facilities that produce them and other manufacturing facilities that use them for industrial processing like the chrome plating industry. Workers who handle PFAS-containing products can also be exposed during their work, such as people who install PFAS-containing carpets and leather furniture with PFAS coatings, professional ski-waxers using PFAS-based waxes, and fire-fighters using PFAS-containing foam and wearing flame-resistant protective gear made with PFASs.
Exposure pathways
People who are exposed to PFASs through their jobs typically have higher levels of PFASs in their blood than the general population. While the general population is exposed to PFASs through ingested food and water, occupational exposure includes accidental ingestion, inhalation exposure, and skin contact in settings where PFAS become volatile.
Professional ski wax technicians
Compared to the general public exposed to contaminated drinking water, professional ski wax technicians are more strongly exposed to PFASs (PFOA, PFNA, PFDA, PFHpA, PFDoDA) from the glide wax used to coat the bottom of skis to reduce the friction between the skis and snow. During the coating process, the wax is heated, which releases fumes and airborne particles. Compared to all other reported occupational and residential exposures, ski waxing had the highest total PFAS air concentrations.
Manufacturing workers
People who work at fluorochemical production plants and in manufacturing industries that use PFASs in the industrial process can be exposed to PFASs in the workplace. Much of what we know about PFAS exposure and health effects began with medical surveillance studies of workers exposed to PFASs at fluorochemical production facilities. These studies began in the 1940s and were conducted primarily at U.S. and European manufacturing sites. Between the 1940s and 2000s, thousands of workers exposed to PFASs participated in research studies that advanced scientific understanding of exposure pathways, toxicokinetic properties, and adverse health effects associated with exposure.
The first research study to report elevated organic fluorine levels in the blood of fluorochemical workers was published in 1980. It established inhalation as a potential route of occupational PFAS exposure by reporting measurable levels of organic fluorine in air samples at the facility. Workers at fluorochemical production facilities have higher levels of PFOA and PFOS in their blood than the general population. Serum PFOA levels in fluorochemical workers are generally below 20,000 ng/mL but have been reported as high as 100,000 ng/mL, whereas the mean PFOA concentration among non-occupationally exposed cohorts in the same time frame was 4.9 ng/mL. Among fluorochemical workers, those with direct contact with PFASs have higher PFAS concentrations in their blood than those with intermittent contact or no direct PFAS contact. Blood PFAS levels have been shown to decline when direct contact ceases. PFOA and PFOS levels have declined in U.S. and European fluorochemical workers due to improved facilities, increased usage of personal protective equipment, and the discontinuation of these chemicals from production. Occupational exposure to PFASs in manufacturing continues to be an active area of study in China with numerous investigations linking worker exposure to various PFASs.
Firefighters
PFASs are commonly used in Class B firefighting foams due to their hydrophobic and lipophobic properties, as well as the stability of the chemicals when exposed to high heat.
Research into occupational exposure for firefighters is emergent, though frequently limited by underpowered study designs. A 2011 cross-sectional analysis of the C8 Health Studies found higher levels of PFHxS in firefighters compared to the sample group of the region, with other PFASs at elevated levels, without reaching statistical significance. A 2014 study in Finland studying eight firefighters over three training sessions observed select PFASs (PFHxS and PFNA) increase in blood samples following each training event. Due to this small sample size, a test of significance was not conducted. A 2015 cross-sectional study conducted in Australia found that PFOS and PFHxS accumulation was positively associated with years of occupational AFFF exposure through firefighting.
Due to their use in training and testing, studies indicate occupational risk for military members and firefighters, as higher levels of PFASs exposure were indicated in military members and firefighters when compared to the general population. PFAS exposure is prevalent among firefighters not only due to its use in emergencies but also because it is used in personal protective equipment. In support of these findings, states like Washington and Colorado have moved to restrict and penalize the use of Class B firefighting foam for firefighter training and testing.
Exposure after September 11 attacks
The September 11 attacks and resulting fires caused the release of toxic chemicals used in materials such as stain-resistant coatings. First responders to this incident were exposed to PFOA, PFNA, and PFHxS through inhalation of dust and smoke released during and after the collapse of the World Trade Center.
Fire responders who were working at or near ground zero were assessed for respiratory and other health effects from exposure to emissions at the World Trade Center. Early clinical testing showed a high prevalence of respiratory health effects. Early symptoms of exposure often presented with persistent coughing and wheezing. PFOA and PFHxS levels were present in both smoke and dust exposure, but first responders exposed to smoke had higher concentrations of PFOA and PFHxS than those exposed to dust.
Mitigation measures
Several strategies have been proposed as a way to protect those who are at greatest risk of occupational exposure to PFAS, including exposure monitoring, regular blood testing, and the use of PFAS-free alternatives such as fluorine-free firefighting foam and plant-based ski wax.
Food and consumer goods
PFASs were found in many plant-based straws, such as paper straws.
Remediation
Water treatment
Several technologies are currently available for remediating PFASs in liquids. These technologies can be applied to drinking water supplies, groundwater, industrial wastewater, surface water, and other applications such as landfill leachate. Influent concentrations of PFASs can vary by orders of magnitude for specific media or applications. These influent values, along with other general water quality parameters (for example, pH) can influence the performance and operating costs of the treatment technologies. The technologies are:
Photodegradation
Foam fractionation
Sorption
Granular activated carbon
Biochar
Ion exchange
Precipitation/flocculation/coagulation
Redox manipulation (chemical oxidation and reduction technologies)
Membrane filtration
Reverse osmosis
Nanofiltration
Supercritical water oxidation
Low Energy Electrochemical Oxidation (EOx)
Private and public sector applications of one or more of these methodologies above are being applied to remediation sites throughout the United States and other international locations. Most solutions involve on-site treatment systems, while others are leveraging off-site infrastructure and facilities, such as a centralized waste treatment facility, to treat and dispose of the PFAS pool of compounds.
The US-based Interstate Technology and Regulatory Council (ITRC) has undertaken an extensive evaluation of ex-situ and in-situ treatment technologies for PFAS-impacted liquid matrices. These technologies are divided into field-implemented technologies, limited application technologies, and developing technologies and typically fit into one of three technology types:
Separation
Concentration
Destruction
The type of PFAS remediation technology selected is often a reflection of the PFAS contamination levels and the PFAS signature (i.e. the combination of short- and long-chain PFAS substances present) in conjunction with the site-specific water chemistry and cross contaminants present in the liquid stream. More complex waters such as landfill leachates and WWTP waters require more robust treatment solutions which are less vulnerable to blockage.
Stripping and enrichment
Foam Fractionation utilizes the air/water interface of a rising air bubble to collect and harvest PFAS molecules. The hydrophobic tail of many long-chain criteria PFAS compounds adhere to this interface and rise to the water surface with the air bubble where they present as a foam for harvesting and further concentration. The foam fractionation technique is a derivation of traditional absorptive bubble separation techniques used by industries for decades to extract amphiphilic contaminants. The absence of a solid absorptive surface reduces consumables and waste byproducts and produces a liquid hyper-concentrate which can be fed into one of the various PFAS destruction technologies. Across various full-scale trials and field applications, this technique provides a simplistic and low operational cost alternative for complex PFAS-impacted waters.
Destruction
In 2007, it was found that high-temperature incineration of sewage sludge reduced the levels of perfluorinated compounds significantly.
A 2022 study published in the Journal of Environmental Engineering found that a heat- and pressure-based technique known as supercritical water oxidation destroyed 99% of the PFAS present in a water sample. During this process, oxidizing substances are added to PFAS-contaminated water and then the liquid is heated above its critical temperature of 374 degrees Celsius at a pressure of more than 220 bars. The water becomes supercritical, and, in this state, water-repellent substances such as PFASs dissolve much more readily.
Theoretical and early-stage solutions
A possible solution for PFAS-contaminated wastewater treatment has been developed by the Michigan State University-Fraunhofer team. Boron-doped diamond electrodes are used for the electrochemical oxidation system where it is capable of breaking PFAS molecular bonds which essentially eliminates the contaminates, leaving fresh water.
Acidimicrobium sp. strain A6 has been shown to be a PFAS and PFOS remediator. PFAS with unsaturated bonds are easier to break down: the commercial dechlorination culture KB1 (contains Dehalococcoides) is capable of breaking down such substances, but not saturated PFAS. When alternative, easier-to-digest substrates are present, microbes may prefer them over PFAS.
Chemical treatment
A study published in Science in August 2022 indicated that perfluoroalkyl carboxylic acids (PFCAs) can be mineralized via heating in a polar aprotic solvent such as dimethyl sulfoxide. Heating PFCAs in an 8 to 1 mixture of dimethyl sulfoxide and water at in the presence of sodium hydroxide caused the removal of the carboxylic acid group at the end of the carbon chain, creating a perfluoroanion that mineralizes into sodium fluoride and other salts such as sodium trifluoroacetate, formate, carbonate, oxalate, and glycolate. The process does not work on perfluorosulfonic acids such as PFOS. A 2022 study published in Chemical Science shows breakdown of C-F bonds and their mineralization as YF3 or YF6 clusters. Another study in the Journal of the American Chemical Society described the PFASs breakdown using metal-organic frameworks (MOFs).
Analytical methods
Analytical methods for PFAS analysis fall into one of two general categories; targeted analysis or non-targeted analysis. Targeted analyses use reference standards to determine concentrations of specific PFAS, but this requires a high-purity standard for each compound of interest. Due to the large number of possible targets, unusual PFAS may go unreported by these methods. Non-targeted analyses measure other factors, such as total organic fluorine, which can be used to estimate the total concentration of PFAS in a sample, but cannot provide concentrations of individual compounds. The two types of analyses are often combined; by subtracting the mass of target analytes from the non-targeted analysis results, one can get an estimate for what fraction of PFAS has been "missed" by the targeted analysis.
Targeted analysis generally use liquid chromatography–mass spectrometry (LC-MS) instruments. Currently, EPA Method 537.1 is approved for use in drinking water and includes 18 PFAS. EPA Method 1633 is undergoing review for use in wastewater, surface water, groundwater, soil, biosolids, sediment, landfill leachate, and fish tissue for 40 PFAS, but is currently being used by many laboratories in the United States. Regulatory limits for PFOA and PFOS set by the US EPA (4 parts-per-trillion) are limited by the capability of methods to detect low level concentrations.
Non-targeted analyses include total organic fluorine (TOF, including variations, e.g., adsorbable organic fluorine, AOF; extractable organic fluorine, EOF), total oxidizable precursor assay, and other methods in development.
Sample chemicals
Some common per- and polyfluoroalkyl substances include:
Polytetrafluoroethylene (aka PTFE or Teflon)
Perfluoroalkyl carboxylic acids (PFCAs)
Perfluorosulfonic acids (PFSAs)
Fluorotelomers (FTOHs)
Litigation and settlements
PFAS have been a subject of multiple lawsuits worldwide. In the United States, settlements stemming from PFAS pollution claims have reached $18 billion by 2024. In 2023, Sweden's Supreme Court set a legal precedent by awarding damages to citizens who were supplied PFAS contaminated drinking water.
In popular culture
Films
The Devil We Know (2018)
Dark Waters (2019)
| Physical sciences | Polymers | Chemistry |
3393371 | https://en.wikipedia.org/wiki/Hindu%E2%80%93Arabic%20numeral%20system | Hindu–Arabic numeral system | The Hindu–Arabic numeral system (also known as the Indo–Arabic numeral system, Hindu numeral system, and Arabic numeral system) is a positional base-ten numeral system for representing integers; its extension to non-integers is the decimal numeral system, which is presently the most common numeral system.
The system was invented between the 1st and 4th centuries by Indian mathematicians. By the 9th century, the system was adopted by Arabic mathematicians who extended it to include fractions. It became more widely known through the writings in Arabic of the Persian mathematician Al-Khwārizmī (On the Calculation with Hindu Numerals, ) and Arab mathematician Al-Kindi (On the Use of the Hindu Numerals, ). The system had spread to medieval Europe by the High Middle Ages, notably following Fibonacci's 13th century Liber Abaci; until the evolution of the printing press in the 15th century, use of the system in Europe was mainly confined to Northern Italy.
It is based upon ten glyphs representing the numbers from zero to nine, and allows representing any natural number by a unique sequence of these glyphs. The symbols (glyphs) used to represent the system are in principle independent of the system itself. The glyphs in actual use are descended from Brahmi numerals and have split into various typographical variants since the Middle Ages.
These symbol sets can be divided into three main families: Western Arabic numerals used in the Greater Maghreb and in Europe; Eastern Arabic numerals used in the Middle East; and the Indian numerals in various scripts used in the Indian subcontinent.
Origins
Sometime around 600 CE, a change began in the writing of dates in the Brāhmī-derived scripts of India and Southeast Asia, transforming from an additive system with separate numerals for numbers of different magnitudes to a positional place-value system with a single set of glyphs for 1–9 and a dot for zero, gradually displacing additive expressions of numerals over the following several centuries.
When this system was adopted and extended by medieval Arabs and Persians, they called it al-ḥisāb al-hindī ("Indian arithmetic"). These numerals were gradually adopted in Europe starting around the 10th century, probably transmitted by Arab merchants; medieval and Renaissance European mathematicians generally recognized them as Indian in origin, however a few influential sources credited them to the Arabs, and they eventually came to be generally known as "Arabic numerals" in Europe. According to some sources, this number system may have originated in Chinese Shang numerals (1200 BCE), which was also a decimal positional numeral system.
Positional notation
The Hindu–Arabic system is designed for positional notation in a decimal system. In a more developed form, positional notation also uses a decimal marker (at first a mark over the ones digit but now more commonly a decimal point or a decimal comma which separates the ones place from the tenths place), and also a symbol for "these digits recur ad infinitum". In modern usage, this latter symbol is usually a vinculum (a horizontal line placed over the repeating digits). In this more developed form, the numeral system can symbolize any rational number using only 13 symbols (the ten digits, decimal marker, vinculum, and a prepended minus sign to indicate a negative number).
Although generally found in text written with the Arabic abjad ("alphabet"), which is written right-to-left, numbers written with these numerals place the most-significant digit to the left, so they read from left to right (though digits are not always said in order from most to least significant). The requisite changes in reading direction are found in text that mixes left-to-right writing systems with right-to-left systems.
Symbols
Various symbol sets are used to represent numbers in the Hindu–Arabic numeral system, most of which developed from the Brahmi numerals.
The symbols used to represent the system have split into various typographical variants since the Middle Ages, arranged in three main groups:
The widespread Western "Arabic numerals" used with the Latin, Cyrillic, and Greek alphabets in the table, descended from the "West Arabic numerals" which were developed in al-Andalus and the Maghreb (there are two typographic styles for rendering western Arabic numerals, known as lining figures and text figures).
The "Arabic–Indic" or "Eastern Arabic numerals" used with Arabic script, developed primarily in what is now Iraq. A variant of the Eastern Arabic numerals is used in Persian and Urdu.
The Indian numerals in use with scripts of the Brahmic family in India and Southeast Asia. Each of the roughly dozen major scripts of India has its own numeral glyphs (as one will note when perusing Unicode character charts).
Glyph comparison
History
Predecessors
The Brahmi numerals at the basis of the system predate the Common Era. They replaced the earlier Kharosthi numerals used since the 4th century BCE. Brahmi and Kharosthi numerals were used alongside one another in the Maurya Empire period, both appearing on the 3rd century BCE edicts of Ashoka.
Buddhist inscriptions from around 300 BCE use the symbols that became 1, 4, and 6. One century later, their use of the symbols that became 2, 4, 6, 7, and 9 was recorded. These Brahmi numerals are the ancestors of the Hindu–Arabic glyphs 1 to 9, but they were not used as a positional system with a zero, and there were rather separate numerals for each of the tens (10, 20, 30, etc.).
The actual numeral system, including positional notation and use of zero, is in principle independent of the glyphs used, and significantly younger than the Brahmi numerals.
Development
The place-value system is used in the Bakhshali manuscript, the earliest leaves being radiocarbon dated to the period 224–383 CE. The development of the positional decimal system Indian mathematics during the Gupta period. Around 500, the astronomer Aryabhata uses the word kha ("emptiness") to mark "zero" in tabular arrangements of digits. The 7th century Brahmasphuta Siddhanta contains a comparatively advanced understanding of the mathematical role of zero. The Sanskrit translation of the lost 5th century Prakrit Jaina cosmological text Lokavibhaga may preserve an early instance of the positional use of zero.
The first dated and undisputed inscription showing the use of a symbol for zero appears on a stone inscription found at the Chaturbhuja Temple at Gwalior in India, dated 876 CE.
Medieval Islamic world
These Indian developments were taken up in Islamic mathematics in the 8th century, as recorded in al-Qifti's Chronology of the scholars (early 13th century).
In 10th century Islamic mathematics, the system was extended to include fractions, as recorded in a treatise by Abbasid Caliphate mathematician Abu'l-Hasan al-Uqlidisi, who was the first to describe positional decimal fractions. According to J. L. Berggren, the Muslims were the first to represent numbers as we do since they were the ones who initially extended this system of numeration to represent parts of the unit by decimal fractions, something that the Hindus did not accomplish. Thus, we refer to the system as "Hindu–Arabic" rather appropriately.
The numeral system came to be known to both the Persian mathematician Khwarizmi, who wrote a book, On the Calculation with Hindu Numerals in about 825 CE, and the Arab mathematician Al-Kindi, who wrote a book, On the Use of the Hindu Numerals ( [kitāb fī isti'māl al-'adād al-hindī]) around 830 CE. Persian scientist Kushyar Gilani who wrote Kitab fi usul hisab al-hind (Principles of Hindu Reckoning) is one of the oldest surviving manuscripts using the Hindu numerals. These books are principally responsible for the diffusion of the Hindu system of numeration throughout the Islamic world and ultimately also to Europe.
Adoption in Europe
In Christian Europe, the first mention and representation of Hindu–Arabic numerals (from one to nine, without zero), is in the (aka Albeldensis), an illuminated compilation of various historical documents from the Visigothic period in Spain, written in the year 976 CE by three monks of the Riojan monastery of San Martín de Albelda. Between 967 and 969 CE, Gerbert of Aurillac discovered and studied Arab science in the Catalan abbeys. Later he obtained from these places the book (On multiplication and division). After becoming Pope Sylvester II in the year 999 CE, he introduced a new model of abacus, the so-called Abacus of Gerbert, by adopting tokens representing Hindu–Arabic numerals, from one to nine.
Leonardo Fibonacci brought this system to Europe. His book introduced Modus Indorum (the method of the Indians), today known as Hindu–Arabic numeral system or base-10 positional notation, the use of zero, and the decimal place system to the Latin world. The numeral system came to be called "Arabic" by the Europeans. It was used in European mathematics from the 12th century, and entered common use from the 15th century to replace Roman numerals.
The familiar shape of the Western Arabic glyphs as now used with the Latin alphabet (0, 1, 2, 3, 4, 5, 6, 7, 8, 9) are the product of the late 15th to early 16th century, when they entered early typesetting. Muslim scientists used the Babylonian numeral system, and merchants used the Abjad numerals, a system similar to the Greek numeral system and the Hebrew numeral system. Similarly, Fibonacci's introduction of the system to Europe was restricted to learned circles. The credit for first establishing widespread understanding and usage of the decimal positional notation among the general population goes to Adam Ries, an author of the German Renaissance, whose 1522 (Calculating on the Lines and with a Quill) was targeted at the apprentices of businessmen and craftsmen.
Adoption in East Asia
In 690 CE, Empress Wu promulgated Zetian characters, one of which was "〇". The word is now used as a synonym for the number zero.
In China, Gautama Siddha introduced Hindu numerals with zero in 718 CE, but Chinese mathematicians did not find them useful, as they had already had the decimal positional counting rods.
In Chinese numerals, a circle (〇) is used to write zero in Suzhou numerals. Many historians think it was imported from Indian numerals by Gautama Siddha in 718 CE, but some Chinese scholars think it was created from the Chinese text space filler "□".
Chinese and Japanese finally adopted the Hindu–Arabic numerals in the 19th century, abandoning counting rods.
Spread of the Western Arabic variant
The "Western Arabic" numerals as they were in common use in Europe since the Baroque period have secondarily found worldwide use together with the Latin alphabet, and even significantly beyond the contemporary spread of the Latin alphabet, intruding into the writing systems in regions where other variants of the Hindu–Arabic numerals had been in use, but also in conjunction with Chinese and Japanese writing (see Chinese numerals, Japanese numerals).
| Mathematics | Basics | null |
2475208 | https://en.wikipedia.org/wiki/Geminal%20diol | Geminal diol | A geminal diol (or gem-diol for short) is any organic compound having two hydroxyl functional groups (-OH) bound to the same carbon atom. Geminal diols are a subclass of the diols, which in turn are a special class of alcohols. Most of the geminal diols are considered unstable.
The simplest geminal diol is methanediol or . Other examples are:
dihydroxymalinic acid
decahydroxycyclopentane
chloral hydrate .
Reactions
Hydration equilibrium
Geminal diols can be viewed as ketone (or aldehyde) hydrates. The two hydroxyl groups in a geminal diol are easily converted to a carbonyl or keto group C=O by loss of one water molecule. Conversely, a keto group can combine with water to form the geminal hydroxyl groups.
The equilibrium in water solution may be shifted towards either compound. For example, the equilibrium constant for the conversion of acetone =O to propane-2,2-diol is about 10−3, while that of formaldehyde =O to methanediol is 103.
For conversion of hexafluoroacetone =O to the diol , the constant is about 10+6, due to the electron withdrawing effect of the trifluoromethyl groups. Similarly, the conversion of chloral =O to chloral hydrate is strongly favored by influence of the trichloromethyl group.
In some cases, such as decahydroxycyclopentane and dodecahydroxycyclohexane, the geminal diol is stable while the corresponding ketone is not.
Geminal diols can also be viewed as extreme cases of hemiacetals, formed by reaction of carbonyl compounds with water, instead of with an alcohol.
| Physical sciences | Alcohols | Chemistry |
2477392 | https://en.wikipedia.org/wiki/Osteolepis | Osteolepis | Osteolepis (from 'bone' and 'scale') is an extinct genus of lobe-finned fish from the Devonian period. It lived in the Lake Orcadie of northern Scotland.
Osteolepis was about long, and covered with large, square scales. The scales and plates on its head were covered in a thin layer of spongy, bony material called cosmine. This layer contained canals that were connected to sensory cells deeper in the skin. These canals ended in pores on the surface and were probably for sensing vibrations in the water.
Osteolepis was a rhipidistian, having a number of features in common with the tetrapods (land-dwelling vertebrates and their descendants), and was probably close to the base of the tetrapod family tree.
| Biology and health sciences | Prehistoric osteichthyans | Animals |
1180073 | https://en.wikipedia.org/wiki/Yamato-class%20battleship | Yamato-class battleship | The were two battleships of the Imperial Japanese Navy, and , laid down leading up to the Second World War and completed as designed. A third hull, laid down in 1940, was converted to the aircraft carrier during construction.
Displacing nearly at full load, the completed battleships were the heaviest ever constructed. The class carried the largest naval artillery ever fitted to a warship, nine 460 mm (18.1 in) naval guns, each capable of firing shells over .
Due to the threat of U.S. submarines and aircraft carriers, both Yamato and Musashi spent the majority of their careers in naval bases at Brunei, Truk, and Kure—deploying on several occasions in response to U.S. raids on Japanese bases.
All three ships were sunk by the U.S. Navy; Musashi by air strikes while participating in the Battle of Leyte Gulf in October 1944, Shinano after being torpedoed by the submarine while under way from Yokosuka to Kure for fitting out in November 1944, and Yamato by air strikes while en route from Japan to Okinawa as part of Operation Ten-Go in April 1945.
Background
The design of the Yamato-class battleships was shaped by expansionist movements within the Japanese government, Japanese industrial power, and the need for a fleet powerful enough to intimidate likely adversaries. Most importantly, the latter, in the form of the Kantai Kessen or Decisive Battle Doctrine, a naval strategy adopted by the Imperial Japanese Navy prior to the Second World War, in which the Japanese navy would win a war by fighting and winning a single, decisive naval action.
After the end of the First World War, many navies—including those of the United States, the United Kingdom, and Imperial Japan—continued and expanded construction programs that had begun during the conflict. The enormous costs associated with these programs pressured their government leaders to begin a disarmament conference. On 8 July 1921, the United States' Secretary of State Charles Evans Hughes invited delegations from the other major maritime powers—France, Italy, Japan, and the United Kingdom—to come to Washington, D.C., and discuss a possible end to the naval arms race. The subsequent Washington Naval Conference resulted in the Washington Naval Treaty. Along with many other provisions, it limited all future battleships to a standard displacement of and a maximum gun caliber of . It also agreed that the five countries would not construct more capital ships for ten years and would not replace any ship that survived the treaty until it was at least twenty years old.
In the 1930s, the Japanese government began a shift towards ultranationalist militancy. This movement called for the expansion of the Japanese Empire to include much of the Pacific Ocean and Southeast Asia. The maintenance of such an empire—spanning from China to Midway Island—required a sizable fleet capable of sustained control of territory. Although all of Japan's battleships built prior to the Yamato class had been completed before 1921—as the Washington Treaty had prevented any more from being completed—all had been either reconstructed or significantly modernized, or both, in the 1930s. This modernization included, among other things, additional speed and firepower, which the Japanese intended to use to conquer and defend their aspired-to empire. When Japan withdrew from the League of Nations in 1934 over the Mukden Incident, it also renounced all treaty obligations, freeing it to build warships larger than those of the other major maritime powers.
Japan's intention to acquire resource-producing colonies in the Pacific and Southeast Asia would likely lead to confrontation with the United States, thus the U.S. became Japan's primary potential enemy. The U.S. possessed significantly greater industrial power than Japan, with 32.2% of worldwide industrial production compared to Japan's 3.5%. Furthermore, several leading members of the United States Congress had pledged "to outbuild Japan three to one in a naval race." Consequently, as Japanese industrial output could not compete with American industrial power, Japanese ship designers developed plans for new battleships individually superior to their counterparts in the United States Navy. Each of these battleships would be capable of engaging multiple enemy capital ships simultaneously, eliminating the need to expend as much industrial effort as the U.S. on battleship construction.
Design
Preliminary studies for a new class of battleships began after Japan's departure from the League of Nations and its renunciation of the Washington and London naval treaties; from 1934 to 1936, 24 initial designs were put forth. These early plans varied greatly in armament, propulsion, endurance, and armor. Main batteries fluctuated between and guns, while the secondary armaments were composed of differing numbers of , , and guns. Propulsion in most of the designs was a hybrid diesel-turbine combination, though one relied solely on diesel and another planned for only turbines. The maximum range of the various designs was between in design A-140-J2 to a high of in designs A-140A and A-140-B2, at a speed of . Armor varied between providing protection from the fire of 406 mm guns to enough protection against 460 mm guns.
After these had been reviewed, two of the original twenty-four were finalized as possibilities, A-140-F3 and A-140-F4. Differing primarily in their range ( versus at ), they were used in the formation of the final preliminary study, which was finished on 20 July 1936. Tweaks to that design resulted in the definitive design of March 1937, which was put forth by Rear-Admiral Fukuda Keiji; a range of 7,200 nmi was finally decided upon, and the hybrid diesel-turbine propulsion was abandoned in favor of turbines. The diesel engines were removed from the design because of problems with the engines aboard the submarine tender Taigei. Their engines, which were similar to the ones that were going to be mounted in the new battleships, required a "major repair and maintenance effort" to keep them running due to a "fundamental design defect". In addition, if the engines failed entirely, the armored citadel deck roof that protected the proposed diesel engine rooms and attendant machinery spaces would severely hamper any attempt to remove and replace them.
The final design called for a standard displacement of and a full-load displacement of , making the ships of the class the largest battleships yet designed, and the largest battleships ever constructed. The design called for a main armament of nine 460 mm naval guns, mounted in three three-gun turrets—each of which weighed more than a 1930s-era destroyer. The designs were quickly approved by the Japanese Naval high command, over the objections of naval aviators, who argued for the construction of aircraft carriers rather than battleships. In all, five Yamato-class battleships were planned.
Ships
Although five Yamato-class vessels had been planned in 1937, only three—two battleships and a converted aircraft carrier—were completed. All three vessels were built in extreme secrecy, to prevent American intelligence officials from learning of their existence and specifications; indeed, the United States' Office of Naval Intelligence only became aware of Yamato and Musashi by name in late 1942. At this early time, their assumptions on the class's specifications were quite far off; while they were correct on their length, the class was given as having a beam of —in actuality, it was about and a displacement of 40,000–57,000 tons (actually, 69,000 tons). In addition, the main armament of Yamato class was given as nine guns as late as July 1945, four months after Yamato was sunk. Both Jane's Fighting Ships and the Western media also misreported the specifications of the ships. In September 1944, Jane's Fighting Ships listed the displacement of both Yamato and Musashi as 45,000 tons. Similarly, both the New York Times and the Associated Press reported that the two ships displaced 45,000 tons with a speed of 30 knots, and even after the sinking of Yamato in April 1945, The Times of London continued to give 45,000 tons as the ship's displacement. Nevertheless, the existence of the ships—and their supposed violation of naval treaties—heavily influenced American naval engineers in the design of the 60,500-ton s, though they were not designed specifically to counter the Yamato class.
Yamato
was ordered in March 1937, laid down 4 November 1937, launched 8 August 1940, and commissioned 16 December 1941. She underwent training exercises until 27 May 1942, when the vessel was deemed "operable" by Admiral Isoroku Yamamoto. Joining the 1st Battleship Division, Yamato served as the flagship of the Japanese Combined Fleet during the Battle of Midway in June 1942, yet did not engage enemy forces during the battle. The next two years were spent intermittently between Truk and Kure naval bases, with her sister ship Musashi replacing Yamato as the flagship of the Combined Fleet. During this time period, Yamato, as part of the 1st Battleship Division, deployed on multiple occasions to counteract American carrier-raids on Japanese island bases. On 25 December 1943, she suffered major torpedo damage at the hands of and was forced to return to Kure for repairs and structural upgrades.
In 1944—following extensive anti-aircraft and secondary battery upgrades—Yamato joined the Second Fleet in the Battle of the Philippine Sea, serving as an escort to a Japanese Carrier Division. In October 1944, as part of Vice-Admiral Takeo Kurita's Center Force for the Battle of Leyte Gulf, she used her naval artillery against an enemy vessel for the only time, helping sink the American escort carrier and the destroyer before she was forced away by torpedoes from , which put her out of combat. Lightly damaged at Kure in March 1945, the ship was then rearmed in preparation for operations. Yamato was deliberately expended in a suicide mission as part of Operation Ten-Go,
sent to use her big guns to provide relief to Japanese forces engaged in the Battle of Okinawa. While en-route she was sunk on 7 April 1945 by 386 American carrier aircraft. After receiving 10 torpedo and 7 bomb hits she capsized, taking 2,498 of the 2,700 crew-members with her, including Vice-Admiral Seiichi Itō. The sinking of Yamato was seen as a major American victory, and Hanson W. Baldwin, the military editor of The New York Times, wrote that "the sinking of the new Japanese battleship Yamato ... is striking proof—if any were needed—of the fatal weakness of Japan in the air and at sea".
Musashi
was ordered in March 1937, laid down 29 March 1938, launched 1 November 1940, and commissioned 5 August 1942. From September to December 1942, she was involved in surface and air-combat training exercises at Hashirajima. On 11 February 1943, Musashi relieved her sister ship Yamato as the flagship of the Combined Fleet. Until July 1944, Musashi shifted between the naval bases of Truk, Yokosuka, Brunei, and Kure. On 29 March 1944, she sustained moderate damage near the bow from one torpedo fired by the American submarine . After repairs and refitting throughout April 1944, Musashi joined the 1st Battleship Division in Okinawa.
In June 1944, as part of the Second Fleet, the ship escorted Japanese aircraft carriers during the Battle of the Philippine Sea. In October 1944, she left Brunei as part of Admiral Takeo Kurita's Center Force during the Battle of Leyte Gulf. Musashi was sunk 24 October during the Battle of the Sibuyan Sea, taking 17 bomb and 19 torpedo hits, with the loss of 1,023 of her 2,399-man crew.
Shinano
Shinano, originally Warship Number 110, was laid down as the third member of the Yamato class, albeit with a slightly modified design. Most of the original armor values were slightly reduced, including the belt, deck, and turrets. The savings in weight this entailed meant that improvements could be made in other areas, including added protection for fire-control and lookout positions. In addition, the secondary armament on the first two Yamatos was to have been replaced by the /65 caliber Type 98 gun. Although smaller, this gun was superior to the 127 mm, possessing a significantly greater muzzle velocity, maximum range, anti-aircraft ceiling, and rate of fire.
In June 1942, following the Japanese defeat at Midway, construction of Shinano was suspended, and the hull was gradually rebuilt as an aircraft carrier. She was designed as a 64,800-ton support vessel that would be capable of ferrying, repairing and replenishing the air fleets of other carriers. Although she was originally scheduled for commissioning in early 1945, the construction of the ship was accelerated after the Battle of the Philippine Sea; this resulted in Shinano being launched on 5 October 1944 and commissioned a little more than a month later on 19 November. Shinano departed Yokosuka for Kure nine days later. In the early morning on 29 November, Shinano was hit by four torpedoes from . Although the damage seemed manageable, poor flooding control caused the vessel to list to starboard. Shortly before midday, she capsized and sank, taking 1,435 of her 2,400-man crew with her. To this day, Shinano is the largest naval vessel to have been sunk by a submarine.
Warships Number 111 and 797
Warship Number 111, never named, was planned as the fourth member of the Yamato class and the second ship to incorporate the improvements of Shinano. The ship's keel was laid after Yamatos launch in August 1940 and construction continued until December 1941, when the Japanese began to question their ambitious capital ship building program—with the coming of war, the resources essential in constructing the ship would become much harder to obtain. As a result, the hull of the fourth vessel, only about 30% complete, was taken apart and scrapped in 1942; materials from this were used in the conversions of and to hybrid battleship/aircraft carriers.
The fifth vessel, Warship Number 797, was planned as an improved Shinano but was never laid down. In addition to the modifications made to that ship, 797 would have removed the two wing turrets in favor of additional 100 mm guns; authors William Garzke and Robert Dulin estimate that this would have allowed for 24 of these weapons. Yamato was eventually modified in 1944 to something akin to this.
Specifications
Armaments
Primary armament
The Yamato-class battleships had primary armaments consisting of three 3-gun turrets mounting /45 caliber Type 94 naval guns – the largest guns ever fitted to a warship, although they were officially designated as the 40 cm/45 caliber (15.9 in) Type 94 – each of which weighed 2,774 tonnes for the complete mount. Each gun was long and weighed , and could fire armor-piercing shells and high explosive shells out to at a rate of 1½ to 2 shells per minute.
The main guns were also capable of firing 3 Shiki tsûjôdan ("Common Type 3") anti-aircraft shells. A time fuze was used to set how far away the shells would explode (although they were commonly set to go off away). Upon detonation, each of these shells would release 900 incendiary-filled tubes in a 20° cone facing towards incoming aircraft; a bursting charge was then used to explode the shell itself to create more steel splinters, finally, the tubes would ignite. The tubes would burn for five seconds at about and would start a flame that was around long. Even though they comprised 40% of the total main ammunition load by 1944, 3 Shiki tsûjôdan were rarely used in combat against enemy aircraft due to the severe damage the firing of these shells inflicted on the barrels of the main guns; indeed, one of the shells may have exploded early and disabled one of Musashis guns during the Battle of the Sibuyan Sea. The shells were intended to put up a barrage of flame that any aircraft attempting to attack would have to navigate through. However, U.S. pilots considered these shells to be more of a pyrotechnics display than a competent anti-aircraft weapon.
Secondary armament
In the original design, the Yamato class' secondary armament comprised twelve 15.5 cm/60 Type 3 guns mounted in four 3-gun turrets (one forward, two amidships, one aft), and twelve 12.7 cm/40 Type 89 guns in six double turrets (three on each side amidships). These had become available once the Mogami-class cruisers were rearmed with guns. With a AP shell, the guns had a maximum range of at an elevation of 45 degrees. Their rate of fire was five rounds per minute. The two midships turrets were removed in 1944 in favor of additional heavy and light anti-aircraft guns.
Initially, heavy anti-aircraft defence was provided by a dozen 40-caliber 127-mm Type 89 dual-purpose guns in six double turrets, three on each side of the superstructure. In 1944, the two amidship 15.5 cm turrets were removed to make room for three additional 127-mm mounts on each side of Yamato, bringing the total number of these guns to twenty-four . When firing at surface targets, the guns had a range of ; they had a maximum ceiling of at their maximum elevation of 90 degrees. Their maximum rate of fire was 14 rounds a minute; their sustained rate of fire was around eight rounds per minute.
Anti-aircraft armament
The Yamato class originally carried twenty-four 25 mm Type 96 anti-aircraft guns, primarily mounted amidships. In 1944, both Yamato and Musashi underwent significant anti-aircraft upgrades in preparation for operations in Leyte Gulf using the space freed up by the removal of both midships secondary battery turrets, and ended up with a complement of twenty-four guns, and one hundred and sixty-two antiaircraft guns, The 25 mm anti-aircraft guns could tilt at 90-degree angles to aim at planes directly overhead, but their mountings' lack of protection made their gunnery crews extremely vulnerable to direct enemy fire. These guns had an effective range of , and an effective ceiling of at an elevation of +85 degrees. The maximum effective rate of fire was only between 110 and 120 rounds per minute because of the frequent need to change the fifteen-round magazines. This was the standard Japanese light AA gun during World War II; it suffered from severe design shortcomings that rendered it a largely ineffective weapon. According to historian Mark Stille, the twin and triple mounts "lacked sufficient speed in train or elevation; the gun sights were unable to handle fast targets; the gun exhibited excessive vibration; the magazine was too small, and ... the gun produced excessive muzzle blast".
The class was also provided with two twin mounts for the licence-built 13.2 mm Type 93 anti-aircraft machine guns, one on each side of the bridge. The maximum range of these guns was , but the effective range against aircraft was only . The cyclic rate was adjustable between 425 and 475 rounds per minute; the need to change 30-round magazines reduced the effective rate to 250 rounds per minute.
The armament on Shinano was quite different from that of her sister vessels due to her conversion. As the carrier was designed for a support role, significant anti-aircraft weaponry was installed on the vessel: sixteen guns, one hundred forty-five anti-aircraft guns, and three hundred and thirty-six anti-aircraft rocket launchers in twelve twenty-eight barrel turrets. None of these guns were ever used against an enemy vessel or aircraft.
Armor
Designed to engage multiple enemy battleships simultaneously, the Yamatos were fitted with heavy armor plating described by naval historian Mark Stille as providing "an unparalleled degree of protection in surface combat". The main belt of armor along the side of the vessel was up to thick, with transverse bulkheads of the armoured citadel up to thick. A lower belt armor thick extending below the main belt was included in the ships as a response to gunnery experiments upon and the new Japanese Type 91 shell which could travel great lengths underwater. Furthermore, the top hull shape was very advanced, the peculiar sideways curving effectively maximizing armor protection and structural rigidity while optimizing weight. The armor on the main turrets surpassed even that of the main belt, with turret face plating thick. Armor plates in both the main belt and main turrets were made of Vickers Hardened steel, which was a face-hardened steel armor. Main armored deck— thick—was composed of a nickel-chromium-molybdenum alloy. Ballistics tests at the proving ground at Kamegakubi demonstrated the deck alloy to be superior to the homogeneous Vickers plates by 10–15%. Additional plating was designed by manipulating the chromium and nickel composition of the alloy. Higher contents of nickel allowed the plate to be rolled and bent without developing fracture properties.
For torpedo protection, a multiple bulkhead side protection system was used which consisted of several void spaces as well as the lower belt armor; the system has a depth of and was designed to withstand a TNT charge. No torpedo defense system compartments were liquid loaded, despite the known benefits. This may have been the result of overestimating the effectiveness of the lower belt armor against torpedoes, an effort to decrease draft, and provision of additional counter-flooding spaces.
The relatively new procedure of arc welding was used extensively throughout the ship, strengthening the durability of the armor plating. Through this technique, the lower-side belt armor was used to strengthen the hull structure of the entire vessel. In total, the vessels of the Yamato class contained 1,147 watertight compartments, of which 1,065 were beneath the armored deck. The ships were also designed with a very large amount of reserve buoyancy to mitigate the effects of flooding.
However, despite the immense armor thickness, the protection scheme of the Yamato class still suffered from several major design flaws and shortcomings. Structural weakness existed near the bow of the vessels, where the armor plating was generally thinner, as demonstrated by Musashi's damage from a torpedo hit in 1943. The hull of the Shinano was subject to even greater structural weakness, being hastily constructed near the end of the war and having been equipped with incomplete armor and unsealed watertight compartments at the time of her sinking. The torpedo defense system performed substantially worse than designed. In particular, very poor jointing between the upper-belt and lower-belt armor created a rupture-prone seam just below the waterline. When combined with the relatively shallow system depth and the lack of liquid loading, this caused the class to be susceptible to torpedoes. Joint failures have been attributed to the considerable damage inflicted upon Yamato from a single torpedo impact in 1943, and to the sinking of Shinano from four hits in 1944.
Propulsion
The Yamato class was fitted with 12 Kampon boilers, which powered quadruple steam turbines, with an indicated horsepower of 147,948 (). These, in turn, drove four propellers. This powerplant enabled the Yamato class to achieve a top speed of . With this speed, the Yamato class' ability to function alongside fast carriers was limited. In addition, the fuel consumption rate of both battleships was very high. As a result, neither battleship was used in combat during the Solomon Islands Campaign or the minor battles during the "island hopping" period of 1943 and early 1944. The propulsion system of Shinano was slightly improved, allowing the carrier to achieve a top speed of .
"Super Yamato"-class battleships
Two battleships of an entirely new and larger design were planned as a part of the 1942 fleet replenishment program. Designated as Design A-150 and initially named Warship Number 178 and Warship Number 179, plans for the ships began soon after the design of the Yamato class was finished, probably in 1938–39. Everything was "essentially completed" sometime in 1941, but with war on the horizon, work on the battleships was halted to fill a need for additional warships, such as aircraft carriers and cruisers, to replace war losses of those vital ships. The Japanese loss in the Battle of Midway, where four carriers were sunk (out of ten, at that time, in the entire navy), made it certain that work on the ships would never begin. In the third volume of their Battleships series, Axis and Neutral Battleships in World War II, the authors William H. Garzke and Robert O. Dulin asserted that these ships would have been the "most powerful battleships in history" because of their massive main battery and extensive anti-aircraft weaponry.
Similar to the fate of papers relating to the Yamato class, most papers and all plans relating to the class were destroyed to prevent capture at the end of the war. It is known that the final design of the ships would have had an even greater firepower and size than the Yamato class—a main battery of six guns in three turrets and secondary dual purpose armament consisting of twenty-four dual mounted guns (similar to the s). The displacement was to be bigger than the Yamatos, and a side armor belt of was planned.
Destruction of records
On the eve of the Allies' occupation of Japan, special-service officers of the Imperial Japanese Navy destroyed virtually all records, drawings, and photographs of or relating to the Yamato-class battleships, leaving only fragmentary records of the design characteristics and other technical matters. The destruction of these documents was so efficient that until 1948 the only known images of Yamato and Musashi were those taken by United States Navy aircraft involved in the attacks on the two battleships. Although some additional photographs and information, from documents that were not destroyed, have come to light over the years, the loss of the majority of written records for the class has made extensive research into the Yamato class somewhat difficult. Because of the lack of written records, information on the class largely came from interviews of Japanese officers following Japan's surrender.
However, in October 1942, based upon a special request from Adolf Hitler, German Admiral Paul Wenneker, attached to the German Naval Attache in Japan, was allowed to inspect a Yamato-class battleship while it was undergoing maintenance in a dockyard, at which time Admiral Wenneker cabled a detailed description of the warship to Berlin. On 22 August 1943, Erich Groner, a German naval historian, and author of the book Die Deutschen Kriegschiffe, 1815–1945, was shown the report while at the "Führer Headquarters", and was directed to make an "interpretation" and then prepare a "design sketch drawing" of the Japanese battleship. The material was preserved by Erich Groner's wife, Mrs. H. Groner, and submitted to publishers in the 1950s.
Cultural significance
From the time of their construction until the present day, Yamato and Musashi have carried a notable presence in Japanese culture, Yamato in particular. Upon completion, the battleships represented the epitome of Imperial Japanese naval engineering. In addition, the two ships, due to their size, speed, and power, visibly embodied Japan's determination and readiness to defend its interests against the western powers, especially the United States. Shigeru Fukudome, chief of the Operations Section of the Imperial Japanese Navy General Staff, described the two ships as "symbols of naval power that provided to officers and men alike a profound sense of confidence in their navy."
Yamato, and especially the story of her sinking, has appeared often in Japanese popular culture, such as the anime Space Battleship Yamato and the 2005 film Yamato. The appearances in popular culture usually portray the ship's last mission as a brave, selfless, but futile, symbolic effort by the participating Japanese sailors to defend their homeland. One of the reasons that the warship may have such significance in Japanese culture is that the word "Yamato" was often used as a poetic name for Japan. Thus, the end of the battleship Yamato could serve as a metaphor for the end of the Japanese empire.
| Technology | Naval warfare | null |
1181008 | https://en.wikipedia.org/wiki/Computational%20science | Computational science | Computational science, also known as scientific computing, technical computing or scientific computation (SC), is a division of science, and more specifically the Computer Sciences, which uses advanced computing capabilities to understand and solve complex physical problems. While this discussion typically extenuates into Visual Computation, this research field of study will typically include the following research categorizations.
Algorithms (numerical and non-numerical): mathematical models, computational models, and computer simulations developed to solve sciences (physical, biological, social), engineering, and humanities problems
Computer hardware that develops and optimizes the advanced system hardware, firmware, networking, and data management components needed to solve computationally demanding problems
The computing infrastructure that supports both the science and engineering problem solving and the developmental computer and information science
In practical use, it is typically the application of computer simulation and other forms of computation from numerical analysis and theoretical computer science to solve problems in various scientific disciplines. The field is different from theory and laboratory experiments, which are the traditional forms of science and engineering. The scientific computing approach is to gain understanding through the analysis of mathematical models implemented on computers. Scientists and engineers develop computer programs and application software that model systems being studied and run these programs with various sets of input parameters. The essence of computational science is the application of numerical algorithms and computational mathematics. In some cases, these models require massive amounts of calculations (usually floating-point) and are often executed on supercomputers or distributed computing platforms.
The computational scientist
The term computational scientist is used to describe someone skilled in scientific computing. Such a person is usually a scientist, an engineer, or an applied mathematician who applies high-performance computing in different ways to advance the state-of-the-art in their respective applied disciplines in physics, chemistry, or engineering.
Computational science is now commonly considered a third mode of science , complementing and adding to experimentation/observation and theory (see image). Here, one defines a system as a potential source of data, an experiment as a process of extracting data from a system by exerting it through its inputs and a model (M) for a system (S) and an experiment (E) as anything to which E can be applied in order to answer questions about S. A computational scientist should be capable of:
recognizing complex problems
adequately conceptualizing the system containing these problems
designing a framework of algorithms suitable for studying this system: the simulation
choosing a suitable computing infrastructure (parallel computing/grid computing/supercomputers)
hereby, maximizing the computational power of the simulation
assessing to what level the output of the simulation resembles the systems: the model is validated
adjusting the conceptualization of the system accordingly
repeat the cycle until a suitable level of validation is obtained: the computational scientist trusts that the simulation generates adequately realistic results for the system under the studied conditions
Substantial effort in computational sciences has been devoted to developing algorithms, efficient implementation in programming languages, and validating computational results. A collection of problems and solutions in computational science can be found in Steeb, Hardy, Hardy, and Stoop (2004).
Philosophers of science addressed the question to what degree computational science qualifies as science, among them Humphreys and Gelfert. They address the general question of epistemology: how does one gain insight from such computational science approaches? Tolk uses these insights to show the epistemological constraints of computer-based simulation research. As computational science uses mathematical models representing the underlying theory in executable form, in essence, they apply modeling (theory building) and simulation (implementation and execution). While simulation and computational science are our most sophisticated way to express our knowledge and understanding, they also come with all constraints and limits already known for computational solutions.
Applications of computational science
Problem domains for computational science/scientific computing include:
Predictive computational science
Predictive computational science is a scientific discipline concerned with the formulation, calibration, numerical solution, and validation of mathematical models designed to predict specific aspects of physical events, given initial and boundary conditions, and a set of characterizing parameters and associated uncertainties. In typical cases, the predictive statement is formulated in terms of probabilities. For example, given a mechanical component and a periodic loading condition, "the probability is (say) 90% that the number of cycles at failure (Nf) will be in the interval N1<Nf<N2".
Urban complex systems
Cities are massively complex systems created by humans, made up of humans, and governed by humans. Trying to predict, understand and somehow shape the development of cities in the future requires complex thinking and computational models and simulations to help mitigate challenges and possible disasters. The focus of research in urban complex systems is, through modeling and simulation, to build a greater understanding of city dynamics and help prepare for the coming urbanization.
Computational finance
In financial markets, huge volumes of interdependent assets are traded by a large number of interacting market participants in different locations and time zones. Their behavior is of unprecedented complexity and the characterization and measurement of the risk inherent to this highly diverse set of instruments is typically based on complicated mathematical and computational models. Solving these models exactly in closed form, even at a single instrument level, is typically not possible, and therefore we have to look for efficient numerical algorithms. This has become even more urgent and complex recently, as the credit crisis has clearly demonstrated the role of cascading effects going from single instruments through portfolios of single institutions to even the interconnected trading network. Understanding this requires a multi-scale and holistic approach where interdependent risk factors such as market, credit, and liquidity risk are modeled simultaneously and at different interconnected scales.
Computational biology
Exciting new developments in biotechnology are now revolutionizing biology and biomedical research. Examples of these techniques are high-throughput sequencing, high-throughput quantitative PCR, intra-cellular imaging, in-situ hybridization of gene expression, three-dimensional imaging techniques like Light Sheet Fluorescence Microscopy, and Optical Projection (micro)-Computer Tomography. Given the massive amounts of complicated data that is generated by these techniques, their meaningful interpretation, and even their storage, form major challenges calling for new approaches. Going beyond current bioinformatics approaches, computational biology needs to develop new methods to discover meaningful patterns in these large data sets. Model-based reconstruction of gene networks can be used to organize the gene expression data in a systematic way and to guide future data collection. A major challenge here is to understand how gene regulation is controlling fundamental biological processes like biomineralization and embryogenesis. The sub-processes like gene regulation, organic molecules interacting with the mineral deposition process, cellular processes, physiology, and other processes at the tissue and environmental levels are linked. Rather than being directed by a central control mechanism, biomineralization and embryogenesis can be viewed as an emergent behavior resulting from a complex system in which several sub-processes on very different temporal and spatial scales (ranging from nanometer and nanoseconds to meters and years) are connected into a multi-scale system. One of the few available options to understand such systems is by developing a multi-scale model of the system.
Complex systems theory
Using information theory, non-equilibrium dynamics, and explicit simulations, computational systems theory tries to uncover the true nature of complex adaptive systems.
Computational science and engineering
Computational science and engineering (CSE) is a relatively new discipline that deals with the development and application of computational models and simulations, often coupled with high-performance computing, to solve complex physical problems arising in engineering analysis and design (computational engineering) as well as natural phenomena (computational science). CSE has become accepted amongst scientists, engineers and academics as the "third mode of discovery" (next to theory and experimentation). In many fields, computer simulation is integral and therefore essential to business and research. Computer simulation provides the capability to enter fields that are either inaccessible to traditional experimentation or where carrying out traditional empirical inquiries is prohibitively expensive. CSE should neither be confused with pure computer science, nor with computer engineering, although a wide domain in the former is used in CSE (e.g., certain algorithms, data structures, parallel programming, high-performance computing), and some problems in the latter can be modeled and solved with CSE methods (as an application area).
Methods and algorithms
Algorithms and mathematical methods used in computational science are varied. Commonly applied methods include:
Computer algebra, including symbolic computation in fields such as statistics, equation solving, algebra, calculus, geometry, linear algebra, tensor analysis (multilinear algebra), optimization
Numerical analysis, including Computing derivatives by finite differences
Application of Taylor series as convergent and asymptotic series
Computing derivatives by Automatic differentiation (AD)
Finite element method for solving PDEs
High order difference approximations via Taylor series and Richardson extrapolation
Methods of integration on a uniform mesh: rectangle rule (also called midpoint rule), trapezoid rule, Simpson's rule
Runge–Kutta methods for solving ordinary differential equations
Newton's method
Discrete Fourier transform
Monte Carlo methods
Numerical linear algebra, including decompositions and eigenvalue algorithms
Linear programming
Branch and cut
Branch and bound
Molecular dynamics, Car–Parrinello molecular dynamics
Space mapping
Time stepping methods for dynamical systems
Historically and today, Fortran remains popular for most applications of scientific computing. Other programming languages and computer algebra systems commonly used for the more mathematical aspects of scientific computing applications include GNU Octave, Haskell, Julia, Maple, Mathematica, MATLAB, Python (with third-party SciPy library), Perl (with third-party PDL library), R, Scilab, and TK Solver. The more computationally intensive aspects of scientific computing will often use some variation of C or Fortran and optimized algebra libraries such as BLAS or LAPACK. In addition, parallel computing is heavily used in scientific computing to find solutions of large problems in a reasonable amount of time. In this framework, the problem is either divided over many cores on a single CPU node (such as with OpenMP), divided over many CPU nodes networked together (such as with MPI), or is run on one or more GPUs (typically using either CUDA or OpenCL).
Computational science application programs often model real-world changing conditions, such as weather, airflow around a plane, automobile body distortions in a crash, the motion of stars in a galaxy, an explosive device, etc. Such programs might create a 'logical mesh' in computer memory where each item corresponds to an area in space and contains information about that space relevant to the model. For example, in weather models, each item might be a square kilometer; with land elevation, current wind direction, humidity, temperature, pressure, etc. The program would calculate the likely next state based on the current state, in simulated time steps, solving differential equations that describe how the system operates, and then repeat the process to calculate the next state.
Conferences and journals
In 2001, the International Conference on Computational Science (ICCS) was first organized. Since then, it has been organized yearly. ICCS is an A-rank conference in the CORE ranking.
The Journal of Computational Science published its first issue in May 2010. The Journal of Open Research Software was launched in 2012.
The ReScience C initiative, which is dedicated to replicating computational results, was started on GitHub in 2015.
Education
At some institutions, a specialization in scientific computation can be earned as a "minor" within another program (which may be at varying levels). However, there are increasingly many bachelor's, master's, and doctoral programs in computational science. The joint degree program master program computational science at the University of Amsterdam and the in computational science was first offered in 2004. In this program, students:
learn to build computational models from real-life observations;
develop skills in turning these models into computational structures and in performing large-scale simulations;
learn theories that will give a firm basis for the analysis of complex systems;
learn to analyze the results of simulations in a virtual laboratory using advanced numerical algorithms.
ETH Zurich offers a bachelor's and master's degree in Computational Science and Engineering. The degree equips students with the ability to understand scientific problem and apply numerical methods to solve such problems. The directions of specializations include Physics, Chemistry, Biology and other Scientific and Engineering disciplines.
George Mason University has offered a multidisciplinary doctorate Ph.D. program in Computational Sciences and Informatics starting from 1992.
The School of Computational and Integrative Sciences, Jawaharlal Nehru University (erstwhile School of Information Technology) also offers a vibrant master's science program for computational science with two specialties: Computational Biology and Complex Systems.
Subfields
Bioinformatics
Car–Parrinello molecular dynamics
Cheminformatics
Chemometrics
Computational archaeology
Computational astrophysics
Computational biology
Computational chemistry
Computational materials science
Computational economics
Computational electromagnetics
Computational engineering
Computational finance
Computational fluid dynamics
Computational forensics
Computational geophysics
Computational history
Computational informatics
Computational intelligence
Computational law
Computational linguistics
Computational mathematics
Computational mechanics
Computational neuroscience
Computational particle physics
Computational physics
Computational sociology
Computational statistics
Computational sustainability
Computer algebra
Computer simulation
Financial modeling
Geographic information science
Geographic information system (GIS)
High-performance computing
Machine learning
Network analysis
Neuroinformatics
Numerical linear algebra
Numerical weather prediction
Pattern recognition
Scientific visualization
Simulation
| Physical sciences | Science basics | Basics and measurement |
1181109 | https://en.wikipedia.org/wiki/Digital%20single-lens%20reflex%20camera | Digital single-lens reflex camera | A digital single-lens reflex camera (digital SLR or DSLR) is a digital camera that combines the optics and mechanisms of a single-lens reflex camera with a solid-state image sensor and digitally records the images from the sensor.
The reflex design scheme is the primary difference between a DSLR and other digital cameras. In the reflex design, light travels through the lens and then to a mirror that alternates to send the image to either a prism, which shows the image in the optical viewfinder, or the image sensor when the shutter release button is pressed. The viewfinder of a DSLR presents an image that will not differ substantially from what is captured by the camera's sensor, as it presents it as a direct optical view through the main camera lens rather than showing an image through a separate secondary lens.
DSLRs largely replaced film-based SLRs during the 2000s. Major camera manufacturers began to transition their product lines away from DSLR cameras to mirrorless interchangeable-lens cameras (MILCs) beginning in the 2010s.
History
In 1969, Willard S. Boyle and George E. Smith invented charge-coupled semiconductor devices, which can be used as analog storage registers and image sensors. A CCD (Charge-Coupled Device) imager provides a low-noise analog image signal, which is digitized when used in a digital camera. For their contribution to digital photography, Boyle and Smith were awarded the Nobel Prize for Physics in 2009.
In 1973, Fairchild developed a 100 x 100 pixel interline CCD image sensor. This CCD was used in the first commercial CCD camera, the Fairchild MV-100, which was introduced in late 1973. In 1974, Kodak scientists Peter Dillon and Albert Brault used this Fairchild CCD 202 image sensor to create the first color CCD image sensor by fabricating a red, green, and blue color filter array that was registered and bonded to the CCD. In 1975, Kodak engineer Steven Sasson built the first portable, battery-operated digital still camera, which used a zoom lens from a Kodak Super 8mm movie camera and a monochrome Fairchild 100×100 pixel CCD.
The first prototype filmless SLR camera was publicly demonstrated by Sony in August 1981. The Sony Mavica (a magnetic still video camera) used a color-striped 2/3” format CCD sensor with 280K pixels, along with analog video signal processing and recording. The Mavica electronic still camera employed a TTL single-lens reflex viewfinder, as shown in the graphic from a June 1982 Sony press release. It recorded FM-modulated analog video signals on a newly developed 2” magnetic floppy disk, dubbed the "Mavipak".
The disk format was later standardized as the "Still Video Floppy", or "SVF", so the Sony Mavica was the first "SVF-SLR" to be demonstrated, but it was not a D-SLR since it recorded analog video images rather than digital images. Starting in 1983, many Japanese companies demonstrated prototype SVF cameras, including Toshiba, Canon, Copal, Hitachi, Panasonic, Sanyo, and Mitsubishi.
The Canon RC-701, introduced in May 1986, was the first SVF camera (and the first SVF-SLR camera) sold in the US. It employed an SLR viewfinder and included a 2/3” format color CCD sensor with 380K pixels. It was sold along with removable 11-66mm and 50-150mm zoom lens.
Over the next five years, many other companies began selling SVF analog electronic cameras. These included the monochrome Nikon QV-1000C SVF-SLR camera, introduced in 1988, which had an F-mount for interchangeable QV Nikkor lenses.
In 1986, the Kodak Microelectronics Technology Division developed a 1.3 MP CCD image sensor, the first with more than 1 million pixels. In 1987, this sensor was integrated with a Canon F-1 film SLR body at the Kodak Federal Systems Division to create an early DSLR camera. The digital back monitored the camera body battery current to sync the image sensor exposure to the film body shutter. Digital images were stored on a tethered hard drive and processed for histogram feedback to the user. This camera was created for the U.S. government, and was followed by several other models intended for government use and eventually Kodak DCS, a commercial DSLR series launched in 1991.
In 1995, Nikon co-developed the Nikon E series with Fujifilm. The E series included the Nikon E2/E2S, Nikon E2N/E2NS and Nikon E3/E3S, with the E3S released in December 1999.
In the late 1990s, Sony introduced the "Digital Mavica" series of consumer digital cameras. Unlike the original analog Mavica, the Digital Mavica cameras recorded JPEG compressed image files on standard 3½-inch magnetic floppy discettes (meant to simplify camera-to-computer data transfer) and did not have an SLR viewfinder.
In 1999, Nikon announced the Nikon D1. The D1's body was similar to Nikon's professional 35 mm film SLRs, and it had the same Nikkor lens mount, allowing the D1 to use Nikon's existing line of AI/AIS manual focus and AF lenses. Although Nikon and other manufacturers had produced digital SLR cameras for several years prior, the D1 was the first professional digital SLR that displaced Kodak's then-undisputed reign over the professional market.
Over the next decade, other camera manufacturers entered the DSLR market, including Canon, Kodak, Fujifilm, Minolta (later Konica Minolta, and ultimately acquired by Sony), Pentax (whose camera division is now owned by Ricoh), Olympus, Panasonic, Samsung, Sigma, and Sony.
In January 2000, Fujifilm announced the FinePix S1 Pro, the first consumer-level DSLR.
In November 2001, Canon released its 4.1-megapixel EOS-1D, the brand's first professional digital body. In 2003, Canon introduced the 6.3-megapixel EOS 300D SLR camera (known in the United States and Canada as the Digital Rebel and in Japan as the Kiss Digital) with an MSRP of US$999, aimed at the consumer market. Its commercial success encouraged other manufacturers to produce competing digital SLRs, lowering entry costs and allowing more amateur photographers to purchase DSLRs.
In 2004, Konica Minolta released the Konica Minolta Maxxum 7D, the first DSLR with in-body image stabilization which later become standard in Pentax, Olympus, and Sony Alpha cameras.
In early 2008, Nikon released the D90, the first DSLR to feature video recording. Since then, all major companies have offered cameras with this functionality.
Over time, the number of megapixels in imaging sensors has increased steadily, with most companies focusing on high ISO performance, speed of focus, higher frame rates, the elimination of digital 'noise' produced by the imaging sensor, and price reductions to lure new customers.
In June 2012, Canon announced the first DSLR to feature a touchscreen, the EOS 650D/Rebel T4i/Kiss X6i. Although this feature had been widely used on both compact cameras and mirrorless models, it had not made an appearance on a DSLR until the 650D.
Market share
The DSLR market is dominated by Japanese companies, and the top five manufacturers are Japanese: Canon, Nikon, Olympus, Pentax, and Sony. Other manufacturers of DSLRs include Mamiya, Sigma, Leica (Germany), and Hasselblad (Swedish).
In 2007, Canon edged out Nikon with 41% of worldwide sales to the latter's 40%, followed by Sony and Olympus, each with approximately 6% market share. In the Japanese domestic market, Nikon captured 43.3% to Canon's 39.9%, with Pentax a distant third at 6.3%.
In 2008, Canon's and Nikon's offerings took the majority of sales. In 2010, Canon controlled 44.5% of the DSLR market, followed by Nikon with 29.8% and Sony with 11.9%.
For Canon and Nikon, digital SLRs are their biggest source of profit. For Canon, their DSLRs brought in four times the profits from compact digital cameras, while Nikon earned more from DSLRs and lenses than from any other product. Olympus and Panasonic have since exited the DSLR market and now focus on producing mirrorless cameras.
In 2013, after a decade of double-digit growth, DSLR (along with MILC) sales were down 15 per cent. This may be due to some low-end DSLR users choosing to use a smartphone instead. The market intelligence firm IDC predicted that Nikon would be out of business by 2018 if the trend continued, although this did not come to pass. Regardless, the market has shifted from being driven by hardware to software, and camera manufacturers have not been keeping up.
Decline and transition to mirrorless cameras
Beginning in the 2010s, major camera manufacturers began to transition their product lines away from DSLR cameras to mirrorless interchangeable-lens cameras (MILCs). In September 2013, Olympus announced they would stop the development of DSLR cameras and focus on the development of MILCs. Nikon announced they were ending production of DSLRs in Japan in 2020, followed by similar announcements from Canon and Sony.
Present-day models
Currently, DSLRs are widely used by consumers and professional still photographers. Well-established DSLRs currently offer a larger variety of dedicated lenses and other equipment. Mainstream DSLRs (in full-frame or smaller image sensor format) are produced by Canon, Nikon, Pentax, and Sigma. Pentax, Phase One, Hasselblad, and Mamiya Leaf produce expensive, high-end medium-format DSLRs, including some with removable sensor backs. Contax, Fujifilm, Kodak, Panasonic, Olympus, Samsung previously produced DSLRs but now either offer non-DSLR systems or have left the camera market entirely. Konica Minolta's line of DSLRs was purchased by Sony.
Canon's current 2018 EOS digital line includes the Canon EOS 1300D/Rebel T6, 200D/SL2, 800D/T7i, 77D, 80D, 7D Mark II, 6D Mark II, 5D Mark IV, 5Ds and 5Ds R and the 1D X Mark II. All Canon DSLRs with three- and four-digit model numbers, as well as the 7D Mark II, have APS-C sensors. The 6D, 5D series, and 1D X are full-frame. , all current Canon DSLRs use CMOS sensors.
Nikon has a broad line of DSLRs, most in direct competition with Canon's offerings, including the D3400, D5600, D7500 and D500 with APS-C sensors, and the D610, D750, D850, D5, D3X and the Df with full-frame sensors.
Leica produces the S2, a medium format DSLR.
Pentax currently offers APS-C, full-frame and medium-format DSLRs. The APS-C cameras include the K-3 II, Pentax KP and K-S2. The K-1 Mark II, announced in 2018 as successor to the Pentax K-1, is the current full-frame model. The APS-C and full-frame models have extensive backward compatibility with Pentax and third-party film era lenses from about 1975, those that use the Pentax K mount. The Pentax 645Z medium format DSLR is also back-compatible with Pentax 645 system lenses from the film era.
Sigma produces DSLRs using the Foveon X3 sensor, rather than the conventional Bayer sensor. This is claimed to give higher colour resolution, although headline pixel counts are lower than conventional Bayer-sensor cameras. It currently offers the entry-level SD15 and the professional SD1. Sigma is the only DSLR manufacturer that sells lenses for other brands' lens mounts.
Sony has modified the DSLR formula in favor of single-lens translucent (SLT) cameras, which are still technically DSLRs, but feature a fixed mirror that allows most light through to the sensor while reflecting some light to the autofocus sensor. Sony's SLTs feature full-time phase detection autofocus during video recording as well as the continuous shooting of up to 12 frame/s. The α series, whether traditional SLRs or SLTs, offers in-body sensor-shift image stabilization and retains the Minolta AF lens mount. , the lineup included the Alpha 68, the semipro Alpha 77 II, and the professional full-frame Alpha 99 II. The translucent (transmissive) fixed mirror allows 70 per cent of the light to pass through onto the imaging sensor, meaning a 1/3rd stop-loss light, but the rest of this light is continuously reflected onto the camera's phase-detection AF sensor for fast autofocus for both the viewfinder and live view on the rear screen, even during the video and continuous shooting. The reduced number of moving parts also makes for faster shooting speeds for its class. This arrangement means that the SLT cameras use an electronic viewfinder as opposed to an optical viewfinder, which some consider a disadvantage, but does have the advantage of a live preview of the shot with current settings, anything displayed on the rear screen is displayed on the viewfinder, and handles bright situations well.
Design
Like SLRs, DSLRs typically use interchangeable lenses with a proprietary lens mount. A movable mechanical mirror system is switched down (to precisely a 45-degree angle) to direct light from the lens over a matte focusing screen via a condenser lens and a pentaprism/pentamirror to an optical viewfinder eyepiece. Most entry-level DSLRs use a pentamirror instead of the traditional pentaprism.
Focusing can be manual, by twisting the focus on the lens; or automatic, activated by pressing half-way on the shutter release or a dedicated auto-focus (AF) button. To take an image, the mirror swings upwards in the direction of the arrow, the focal-plane shutter opens, and the image is projected and captured on the image sensor. After these actions, the shutter closes, the mirror returns to the 45-degree angle, and the built-in drive mechanism re-tensions the shutter for the next exposure.
Compared with the newer concept of mirrorless interchangeable-lens cameras, this mirror/prism system is the characteristic difference, providing direct, accurate optical preview with separate autofocus and exposure metering sensors. Essential parts of all digital cameras are some electronics like amplifiers, analog-to-digital converters, image processors, and other microprocessors for processing the digital image, performing data storage, and/or driving an electronic display.
DSLRs typically use autofocus based on phase detection. This method allows the optimal lens position to be calculated rather than "found", as would be the case with autofocus based on contrast maximization. Phase-detection autofocus is typically faster than other passive techniques. As the phase sensor requires the same light going to the image sensor, it was previously only possible with an SLR design. However, with the introduction of focal-plane phase-detect autofocusing in mirrorless interchangeable lens cameras by Sony, Fuji, Olympus, and Panasonic, cameras can now employ both phase detect and contrast-detect AF points.
Common features
Mode dial
Digital SLR cameras, along with most other digital cameras, generally have a mode dial to access standard camera settings or automatic scene-mode settings. Sometimes called a "PASM" dial, they typically provide modes such as program, aperture-priority, shutter-priority, and full manual modes. Scene modes vary from camera to camera, and these modes are inherently less customizable. They often include landscape, portrait, action, macro, night, and silhouette, among others. However, these different settings and shooting styles that "scene" mode provides can be achieved by calibrating certain settings on the camera.
Dust reduction systems
A method to prevent dust from entering the chamber by using a "dust cover" filter right behind the lens mount was used by Sigma in its first DSLR, the Sigma SD9, in 2002.
Olympus used a built-in sensor cleaning mechanism in its first DSLR that had a sensor exposed to air, the Olympus E-1, in 2003 (all previous models each had a non-interchangeable lens, preventing direct exposure of the sensor to outside environmental conditions).
Several Canon DSLR cameras rely on dust reduction systems based on vibrating the sensor at ultrasonic frequencies to remove dust from the sensor.
Interchangeable lenses
The ability to exchange lenses, to select the best lens for the current photographic need, and to allow the attachment of specialized lenses is one of the key factors in the popularity of DSLR cameras, although this feature is not unique to the DSLR design and mirrorless interchangeable lens cameras are becoming increasingly popular. Interchangeable lenses for SLRs and DSLRs are built to operate correctly with a specific lens mount that is generally unique to each brand. A photographer will often use lenses made by the same manufacturer as the camera body (for example, Canon EF lenses on a Canon body) although there are also many independent lens manufacturers, such as Sigma, Tamron, Tokina, and Vivitar, that make lenses for a variety of different lens mounts. There are also lens adapters that allow a lens for one lens mount to be used on a camera body with a different lens mount, but with often reduced functionality.
Many lenses are mountable, "diaphragm-and-meter-compatible", on modern DSLRs, and on older film SLRs that use the same lens mount. However, when lenses designed for 35 mm film or equivalently sized digital image sensors are used on DSLRs with smaller sized sensors, the image is effectively cropped and the lens appears to have a longer focal length than its stated focal length. Most DSLR manufacturers have introduced lines of lenses with image circles optimised for the smaller sensors and focal lengths equivalent to those generally offered for existing 35 mm mount DSLRs, mostly in the wide-angle range. These lenses tend not to be completely compatible with full-frame sensors or 35 mm film because of the smaller imaging circle and with some Canon EF-S lenses, interfere with the reflex mirrors on full-frame bodies.
HD video capture
Since 2008, manufacturers have offered DSLRs which offer a movie mode capable of recording high definition motion video. A DSLR with this feature is often known as an HDSLR or DSLR video shooter. The first DSLR introduced with an HD movie mode, the Nikon D90, captures video at 720p24 (1280x720 resolution at 24 frame/s). Other early HDSLRs capture video using a nonstandard video resolution or frame rate. For example, the Pentax K-7 uses a nonstandard resolution of 1536×1024, which matches the imager's 3:2 aspect ratio. The Canon EOS 500D (Rebel T1i) uses a nonstandard frame rate of 20 frame/s at 1080p, along with a more conventional 720p30 format.
In general, HDSLRs use the full imager area to capture HD video, though not all pixels (causing video artifacts to some degree). Compared with the much smaller image sensors found in the typical camcorder, the HDSLR's much larger sensor yields distinctly different image characteristics. HDSLRs can achieve much shallower depth of field and superior low-light performance. However, the low ratio of active pixels (to total pixels) is more susceptible to aliasing artifacts (such as moiré patterns) in scenes with particular textures, and CMOS rolling shutter tends to be more severe. Furthermore, due to the DSLR's optical construction, HDSLRs typically lack one or more video functions found on standard dedicated camcorders, such as autofocus while shooting, powered zoom, and an electronic viewfinder/preview. These and other handling limitations prevent the HDSLR from being operated as a simple point-and-shoot camcorder, instead of demanding some level of planning and skill for location shooting.
Video functionality has continued to improve since the introduction of the HDSLR, including higher video resolution (such as 1080p24) and video bitrate, improved automatic control (autofocus) and manual exposure control, and support for formats compatible with high-definition television broadcast, Blu-ray disc mastering or Digital Cinema Initiatives (DCI). The Canon EOS 5D Mark II (with the release of firmware version 2.0.3/2.0.4.) and Panasonic Lumix GH1 were the first HDSLRs to offer 1080p video at 24fps, and since then the list of models with comparable functionality has grown considerably.
The rapid maturation of HDSLR cameras has sparked a revolution in digital filmmaking (referred to as "DSLR revolution"), and the "Shot On DSLR" badge is a quickly growing phrase among independent filmmakers. Canon's North American TV advertisements featuring the Rebel T1i have been shot using the T1i itself. Other types of HDSLRs found their distinct application in the field of documentary and ethnographic filmmaking, especially due to their affordability, technical and aesthetical features, and their ability to make observation highly intimate. An increased number of films, television shows, and other productions are utilizing the quickly improving features. One such project was Canon's "Story Beyond the Still" contest that asked filmmakers to collectively shoot a short film in 8 chapters, with each chapter being shot over a short period of time and a winner was determined for each chapter. After 7 chapters the winners collaborated to shoot the final chapter of the story. Due to the affordability and convenient size of HDSLRs compared with professional movie cameras, The Avengers used five Canon EOS 5D Mark II and two Canon 7D to shoot the scenes from various vantage angles throughout the set and reduced the number of reshoots of complex action scenes.
Manufacturers have sold optional accessories to optimize a DSLR camera as a video camera, such as a shotgun-type microphone, and an External EVF with 1.2 million pixels.
Live preview
Early DSLRs lacked the ability to show the optical viewfinder's image on the LC display – a feature known as live preview. Live preview is useful in situations where the camera's eye-level viewfinder cannot be used, such as underwater photography where the camera is enclosed in a plastic waterproof case.
In 2000, Olympus introduced the Olympus E-10, the first DSLR with live preview – albeit with an atypical fixed lens design. , some DSLRs from Canon, Nikon, Olympus, Panasonic, Leica, Pentax, Samsung and Sony all provided continuous live preview as an option. Additionally, the Fujifilm FinePix S5 Pro offers 30 seconds of live preview.
On almost all DSLRs that offer live preview via the primary sensor, the phase-detection autofocus system does not work in the live preview mode, and the DSLR switches to a slower contrast system commonly found in point-and-shoot cameras. While even phase detection autofocus requires contrast in the scene, strict contrast-detection autofocus is limited in its ability to find focus quickly, though it is somewhat more accurate.
In 2012, Canon introduced hybrid autofocus technology to the DSLR in the EOS 650D/Rebel T4i, and introduced a more sophisticated version, which it calls "Dual Pixel CMOS AF", with the EOS 70D. The technology allows certain pixels to act as both contrast-detection and phase-detection pixels, thereby greatly improving autofocus speed in live view (although it remains slower than pure phase detection). While several mirrorless cameras, plus Sony's fixed-mirror SLTs, have similar hybrid AF systems, Canon is the only manufacturer that offers such technology in DSLRs.
A new feature via a separate software package introduced from Breeze Systems in October 2007, features live view from a distance. The software package is named "DSLR Remote Pro v1.5" and enables support for the Canon EOS 40D and 1D Mark III.
Sensor size and image quality
Image sensors used in DSLRs come in a range of sizes. The very largest are the ones used in "medium format" cameras, typically via a "digital back" which can be used as an alternative to a film back. Because of the manufacturing costs of these large sensors, the price of these cameras is typically over $1,500 and easily reaching $8,000 and beyond .
"Full-frame" is the same size as 35 mm film (135 film, image format 24×36 mm); these sensors are used in DSLRs such as the Canon EOS-1D X Mark II, 5DS/5DSR, 5D Mark IV and 6D Mark II, and the Nikon D5, D850, D750, D610 and Df. Most lower-cost DSLRs use a smaller sensor that is APS-C sized, which is approximately 24×16 mm, slightly smaller than the size of an APS-C film frame, or about 40% of the area of a full-frame sensor. Other sensor sizes found in DSLRs include the Four Thirds System sensor at 26% of full frame, APS-H sensors (used, for example, in the Canon EOS-1D Mark III) at around 61% of full frame, and the original Foveon X3 sensor at 33% of full frame (although Foveon sensors since 2013 have been APS-C sized). Leica offers an "S-System" DSLR with a 30×45 mm array containing 37 million pixels. This sensor is 56% larger than a full-frame sensor.
The resolution of DSLR sensors is typically measured in megapixels. More expensive cameras and cameras with larger sensors tend to have higher megapixel ratings. A larger megapixel rating does not mean higher quality. Low light sensitivity is a good example of this. When comparing two sensors of the same size, for example, two APS-C sensors one 12.1 MP and one 18 MP, the one with the lower megapixel rating will usually perform better in low light. This is because the size of the individual pixels is larger, and more light is landing on each pixel, compared with the sensor with more megapixels. This is not always the case, because newer cameras that have higher megapixels also have better noise reduction software, and higher ISO settings to make up for the loss of light per pixel due to higher pixel density.
Depth-of-field control
The lenses typically used on DSLRs have a wider range of apertures available to them, ranging from as large as 0.9 to about 32. Lenses for smaller sensor cameras rarely have true available aperture sizes much larger than 2.8 or much smaller than 5.6.
To help extend the exposure range, some smaller sensor cameras will also incorporate an ND filter pack into the aperture mechanism.
The apertures that smaller sensor cameras have available give much more depth of field than equivalent angles of view on a DSLR. For example, a 6 mm lens on a 2/3″ sensor digicam has a field of view similar to a 24 mm lens on a 35 mm camera. At an aperture of 2.8, the smaller sensor camera (assuming a crop factor of 4) has a similar depth of field to that 35 mm camera set to 11.
Wider angle of view
The angle of view of a lens depends upon its focal length and the camera's image sensor size; a sensor smaller than 35 mm film format (36×24 mm frame) gives a narrower angle of view for a lens of a given focal length than a camera equipped with a full-frame (35 mm) sensor. As of 2017, only a few current DSLRs have full-frame sensors, including the Canon EOS-1D X Mark II, EOS 5D Mark IV, EOS 5DS/5DS R, and EOS 6D Mark II; Nikon's D5, D610, D750, D850, and Df; and the Pentax K-1. The scarcity of full-frame DSLRs is partly a result of the cost of such large sensors. Medium format size sensors, such as those used in the Mamiya ZD among others, are even larger than full-frame (35 mm) sensors, and capable of even greater resolution, and are correspondingly more expensive.
The impact of sensor size on the field of view is referred to as the "crop factor" or "focal length multiplier", which is a factor by which a lens focal length can be multiplied to give the full-frame-equivalent focal length for a lens. Typical APS-C sensors have crop factors of 1.5 to 1.7, so a lens with a focal length of 50 mm will give a field of view equal to that of a 75 mm to 85 mm lens on a 35 mm camera. The smaller sensors of Four Thirds System cameras have a crop factor of 2.0.
While the crop factor of APS-C cameras effectively narrows the angle of view of long-focus (telephoto) lenses, making it easier to take close-up images of distant objects, wide-angle lenses suffer a reduction in their angle of view by the same factor.
DSLRs with "crop" sensor size have slightly more depth-of-field than cameras with 35 mm sized sensors for a given angle of view. The amount of added depth of field for a given focal length can be roughly calculated by multiplying the depth of field by the crop factor. Shallower depth of field is often preferred by professionals for portrait work and to isolate a subject from its background.
Unusual features
On July 13, 2007, FujiFilm announced the FinePix IS Pro, which uses Nikon F-mount lenses. This camera, in addition to having live preview, has the ability to record in the infrared and ultraviolet spectra of light.
In August 2010 Sony released series of DSLRs allowing 3D photography. It was accomplished by sweeping the camera horizontally or vertically in Sweep Panorama 3D mode. The picture could be saved as ultra-wide panoramic image or as 16:9 3D photography to be viewed on BRAVIA 3D television set.
Comparison with other digital cameras
The reflex design scheme is the primary difference between a DSLR and other digital cameras. In the reflex design scheme, the image captured on the camera's sensor is also the image that is seen through the viewfinder. Light travels through a single lens and a mirror is used to reflect a portion of that light through the viewfinder – hence the name "single-lens reflex". While there are variations among point-and-shoot cameras, the typical design exposes the sensor constantly to the light projected by the lens, allowing the camera's screen to be used as an electronic viewfinder. However, LCDs can be difficult to see in very bright sunlight.
Compared with some low-cost cameras that provide an optical viewfinder that uses a small auxiliary lens, the DSLR design has the advantage of being parallax-free: it never provides an off-axis view. A disadvantage of the DSLR optical viewfinder system is that when it is used, it prevents using the LCD for viewing and composing the picture. Some people prefer to compose pictures on the display – for them, this has become the de facto way to use a camera. Depending on the viewing position of the reflex mirror (down or up), the light from the scene can only reach either the viewfinder or the sensor. Therefore, many early DSLRs did not provide "live preview" (i.e., focusing, framing, and depth-of-field preview using the display), a facility that is always available on digicams. Today most DSLRs can alternate between live view and viewing through an optical viewfinder.
Optical view image and digitally created image
The larger, advanced digital cameras offer a non-optical electronic through-the-lens (TTL) view, via an eye-level electronic viewfinder (EVF) in addition to the rear LCD. The difference in view compared with a DSLR is that the EVF shows a digitally created image, whereas the viewfinder in a DSLR shows an actual optical image via the reflex viewing system. An EVF image has the lag time (that is, it reacts with a delay to view changes) and has a lower resolution than an optical viewfinder but achieves parallax-free viewing using less bulk and mechanical complexity than a DSLR with its reflex viewing system. Optical viewfinders tend to be more comfortable and efficient, especially for action photography and in low-light conditions. Compared with digital cameras with LCD electronic viewfinders, there is no time lag in the image: it is always correct as it is being "updated" at the speed of light. This is important for action or sports photography, or any other situation where the subject or the camera is moving quickly. Furthermore, the "resolution" of the viewed image is much better than that provided by an LCD or an electronic viewfinder, which can be important if manual focusing is desired for precise focusing, as would be the case in macro photography and "micro-photography" (with a microscope). An optical viewfinder may also cause less eye-strain. However, electronic viewfinders may provide a brighter display in low light situations, as the picture can be electronically amplified.
Performance differences
DSLR cameras often have image sensors of much larger size and often higher quality, offering lower noise, which is useful in low light. Although mirrorless digital cameras with APS-C and full frame sensors exist, most full frame and medium format sized image sensors are still seen in DSLR designs.
For a long time, DSLRs offered faster and more responsive performance, with less shutter lag, faster autofocus systems, and higher frame rates. Around 2016–17, some mirrorless camera models started offering competitive or superior specifications in these aspects. The downside of these cameras being that they do not have an optical viewfinder, making it difficult to focus on moving subjects or in situations where a fast burst mode would be beneficial. Other digital cameras were once significantly slower in image capture (time measured from pressing the shutter release to the writing of the digital image to the storage medium) than DSLR cameras, but this situation is changing with the introduction of faster capture memory cards and faster in-camera processing chips. Still, compact digital cameras are not suited for action, wildlife, sports, and other photography requiring a high burst rate (frames per second).
Simple point-and-shoot cameras rely almost exclusively on their built-in automation and machine intelligence for capturing images under a variety of situations and offer no manual control over their functions, a trait that makes them unsuitable for use by professionals, enthusiasts, and proficient consumers (also known as "prosumers"). Bridge cameras provide some degree of manual control over the camera's shooting modes, and some even have hot shoes and the option to attach lens accessories such as filters and secondary converters. DSLRs typically provide the photographer with full control over all the important parameters of photography and have the option to attach additional accessories using the hot shoe. including hot shoe-mounted flash units, battery grips for additional power and hand positions, external light meters, and remote controls. DSLRs typically also have fully automatic shooting modes.
DSLRs have a larger focal length for the same field of view, which allows the creative use of depth of field effects. However, small digital cameras can focus better on closer objects than typical DSLR lenses.
The sensors used in current DSLRs — "full-frame" which is the same size as 35mm film, APS-C, and Four Thirds System — are much larger than most digital cameras. Entry-level compact cameras typically use sensors known as 1/2.3″, which is 3% the size of a full-frame sensor. There are fixed-lens cameras — such as bridge cameras, premium compact cameras, or high-end point-and-shoot cameras — that offer sensors larger than 1/2.3″, but many still fall short of the larger sizes widely found in DSLRs. Examples include the Sigma DP1, which uses a Foveon X3 sensor; the Leica X1; the Canon PowerShot G1 X, which uses a 1.5″ (18.7×14 mm) sensor that is slightly larger than the Four Thirds standard and is 30% of a full-frame sensor; the Nikon Coolpix A, which uses an APS-C sensor of the same size as those found in the company's DX-format DSLRs; and two models from Sony, the RX100 with a 1″-type (13.2×8.8 mm) sensor with about half the area of Four Thirds and the full-frame Sony RX1. These premium compacts are often comparable to entry-level DSLRs in price, with a smaller sensor being a tradeoff for the size and weight savings.
Fixed or interchangeable lenses
Unlike DSLRs, most digital cameras lack the option to change the lens. Instead, most compact digital cameras are manufactured with a zoom lens that covers the most commonly used fields of view. Having fixed lenses, they are limited to the focal lengths they are manufactured with, except for what is available from attachments. Manufacturers have attempted (with increasing success) to overcome this disadvantage by offering extreme ranges of focal length on models known as superzooms, some of which offer far longer focal lengths than readily available DSLR lenses.
There are now available perspective-correcting (PC) lenses for DSLR cameras, providing some of the attributes of view cameras. Nikon introduced the first fully manual PC lens in 1961. Recently, however, some manufacturers have introduced advanced lenses that shift and tilt and are operated with automatic aperture control.
However, since the introduction of the Micro Four Thirds system by Olympus and Panasonic in late 2008, mirrorless interchangeable lens cameras are now widely available. Hence, the option to change lenses is no longer unique to DSLRs. Cameras for the micro four-thirds system are designed with the option of a replaceable lens, and lenses that conform to this proprietary specification are accepted. Cameras for this system have the same sensor size as the Four-Thirds System but do not have the mirror and pentaprism to reduce the distance between the lens and sensor.
Panasonic released the first Micro Four Thirds camera, the Lumix DMC-G1. Several manufacturers have announced lenses for the new Micro Four Thirds mount. In contrast, older Four-Thirds lenses can be mounted with an adapter (a mechanical spacer with front and rear electrical connectors and its own internal firmware). A similar mirror-less interchangeable lens camera with an APS-C-sized sensor was announced in January 2010: the Samsung NX10. On 21 September 2011, Nikon announced with the Nikon 1 a series of high-speed MILCs. A handful of rangefinder cameras also support interchangeable lenses. Six digital rangefinders exist: the Epson R-D1 (APS-C-sized sensor), the Leica M8 (APS-H-sized sensor), both smaller than 35 mm film rangefinder cameras, and the Leica M9, M9-P, M Monochrom and M (Typ 240) (all full-frame cameras, with the Monochrom shooting exclusively in black-and-white).
In common with other interchangeable lens designs, DSLRs must contend with potential sensor contamination by dust particles when the lens is changed (though recent dust reduction systems alleviate this). Digital cameras with fixed lenses are not usually subject to dust from outside the camera settling on the sensor.
DSLRs generally have more significant cost, size, and weight. They also have louder operation, due to the SLR mirror mechanism. Sony's fixed mirror design manages to avoid this problem. However, that design has the disadvantage that the mirror diverts some of the light received from the lens, and thus, the image sensor receives about 30% less light compared with other DSLR designs.
| Technology | Photography | null |
1182474 | https://en.wikipedia.org/wiki/Viridian | Viridian | Viridian is a blue-green pigment, a hydrated chromium(III) oxide, of medium saturation and relatively dark in value. It is composed of a majority of green, followed by blue. The first recorded use of viridian as a color name in English was in the 1860s. Viridian takes its name from the Latin viridis, meaning "green". The pigment was first prepared in mid-19th-century Paris and remains available from several US manufacturers as prepared artists' colors in all media.
History
Viridian pigment was first prepared in 1838 in Paris by Parisian color chemist and painter Pannetier alongside his assistant Binet as a hydrated form of chromium oxide. The preparation process was demanding, expensive, and shrouded in secrecy. The French chemist C. E. Guignet developed and patented a cheaper manufacturing method in 1859 that enabled larger distribution and use of the pigment. This method involved calcining a combination of boric acid and potassium bichromate, then washing the material.
Winsor and Newton's catalogue listed the pigment as early as 1849. It was used as early as 1840 in a work by J. M. W. Turner. Viridian was in prominent use by the mid-nineteenth century, but was less popular than three to four times more affordable alternatives including emerald and chrome greens.
Visual characteristics
Viridian is a bright shade of spring green, which places the color between green and teal on the color wheel, or, in paint, a tertiary blue–green color. Viridian is dark in value, has medium saturation, and is transparent .
Variations of viridian
Paolo Veronese green
Paolo Veronese green is the color that is called Verde Verones in the Guía de coloraciones (Guide to colorations) by Rosa Gallego and Juan Carlos Sanz, a color dictionary published in 2005 that is widely popular in the Hispanophone realm.
Paolo Veronese green was a color formulated and used by the noted 16th-century Venetian artist Paolo Veronese.
Paolo Veronese green began to be used as a color name in English sometime in the 1800s (exact year uncertain).
Another name for this color is transparent oxide of chromium.
Viridian green
At right is displayed the color viridian green.
The source of this color is the "Pantone Textile Paper eXtended (TPX)" color list, color #17-5126 TPX—Viridian Green.
Generic viridian
Generic viridian is the color that is called Viridian inspecifico in the Guía de coloraciones (Guide to colorations) by Rosa Gallego and
Juan Carlos Sanz, a color dictionary published in 2005 that is widely popular in the Hispanophone realm.
Spanish viridian
Spanish viridian is the color that is called Viridian specifico in the Guía de coloraciones (Guide to colorations) by Rosa Gallego and
Juan Carlos Sanz, a color dictionary published in 2005 that is widely popular in the Hispanophone realm.
Permanence
Viridian is considered durable and permanent as an artist's pigment. Viridian is unaffected by temperatures up to 260 °C (500 °F), but it is unsuitable for use in ceramic glazes. Viridian is compatible with all pigments in all media, and has high oil absorption. Pure pigment formulations of viridian are hard and may separate in tubes, but adding barium sulfate in small quantities enables easy grinding and dispersion.
Notable occurrences
Although viridian is not a frequent color name in English, it is used in a number of cultural references, probably because it is derived from viridis, the Latin word for green, so using the word viridian sounds more elegant than simply referring to the Old English word green.
Fine art painting
Fritz Bamberger, Afterglow in the Sierra Nevada, 1863.
Claude Monet, Arrival of the Normandy Train, Gare Saint-Lazare, 1877, oil on canvas includes traces of viridian in the grassy area.
Pierre-Auguste Renoir, Flowers, 1919.
Automobiles
"Viridian Joule" was the winning color name in Chevrolet's Volt Paint-Color Naming Contest.
Broadcasting
Viridian was the signature color of BBC Two's identity from 1991–2001.
Television
Viridian was the name of an Orion starship in Star Trek: Discovery
Environmental design
The viridian design movement is a popular design movement based on a bright green environmentalism philosophy.
Film
Viridian is mentioned by Otho when discussing remodeling, in the 1988 film Beetlejuice.
Music
"Viridian" is the seventh song on Between the Buried and Me's Colors.
"Viridian (Interlude)" is the sixteenth track on Bethel Music's Without Words: Synesthesia
Viridian is the name of the ninth studio album by rock band Closterkeller whose releases are usually named after colors
"Viridian" is the fourth song in Novo Amor's 2022 instrumental album Antarctican Dream Machine
Literature
In the Space opera series Dread Empire's Fall, by American author Walter Jon Williams, the first mention of viridian is in the initial book of the series, The Praxis for both the sky color of Shaa Empire capital world Zanshaa as well as color of the uniforms of their Fleet.
Video games
In the Pokémon franchise, in the Kanto region, Viridian City is the first town one encounters after leaving Pallet Town via Route 1 and also home to the final gym.
In VVVVVV, the player character is Captain Viridian, who is a light blue-green color. All characters have names referencing their color and starting with the letter V.
In Knights of the Old Republic II the player character can find and use a viridian lightsaber crystal.
In League of Legends the champion Kayle has a viridian costume that is green with black wings.
In Phoenix Wright: Ace Attorney − Trials and Tribulations, there is a running joke about the color, started by the character Larry Butz.
In the Steam game Aviary Attorney, the Viridian Killer is responsible for murders in France during the 1830s.
| Physical sciences | Colors | Physics |
1182894 | https://en.wikipedia.org/wiki/Turquoise%20%28color%29 | Turquoise (color) | Turquoise ( ) is a cyan color, based on the mineral of the same name. The word turquoise dates to the 17th century and is derived from the French , meaning 'Turkish', because the mineral was first brought to Europe through Turkey from mines in the historical Khorasan province of Iran (Persia) and Afghanistan today. The first recorded use of turquoise as a color name in English was in 1573.
The X11 color named turquoise is displayed on the right.
Turquoise gemstones
Turquoise is an opaque, blue-to-green mineral that is a hydrous phosphate of copper and aluminium, with the chemical formula CuAl6(PO4)4(OH)8·4H2O. It is rare and valuable in finer grades and has been prized as a gem and ornamental stone for thousands of years owing to its unique hue.
In many cultures of the Old and New Worlds, this gemstone has been esteemed for thousands of years as a holy stone, a bringer of good fortune or a talisman. The oldest evidence for this claim was found in ancient Egypt, where grave furnishings with turquoise inlay were discovered, dating from approximately 3000 BCE. In the ancient Persian Empire, the sky-blue gemstones were earlier worn round the neck or wrist as protection against unnatural death. If they changed color, the wearer was thought to have reason to fear the approach of doom. Meanwhile, it has been discovered that turquoise can change color. The change can be caused by light, or by a chemical reaction brought about by cosmetics, dust, or the acidity of the skin.
Turquoise is a stone and color that is strongly associated with the domes and interiors of large mosques in Iran, Central Asia, and Russia.
Variations
Celeste
Celeste is a sky blue turquoise.
Light turquoise
Light turquoise is a lighter tone of turquoise.
Turquoise blue
Turquoise blue is close to turquoise on the color wheel, but slightly more blue.
The first recorded use of turquoise blue as a color name in English was in 1900.
Medium turquoise
This is the web color medium turquoise.
Dark turquoise
This is the web color dark turquoise.
Bright turquoise
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| Physical sciences | Colors | Physics |
1182980 | https://en.wikipedia.org/wiki/Mountain%20formation | Mountain formation | Mountain formation occurs due to a variety of geological processes associated with large-scale movements of the Earth's crust (tectonic plates). Folding, faulting, volcanic activity, igneous intrusion and metamorphism can all be parts of the orogenic process of mountain building. The formation of mountains is not necessarily related to the geological structures found on it.
From the late 18th century until its replacement by plate tectonics in the 1960s, geosyncline theory was used to explain much mountain-building. The understanding of specific landscape features in terms of the underlying tectonic processes is called tectonic geomorphology, and the study of geologically young or ongoing processes is called neotectonics.
Types of mountains
There are five main types of mountains: volcanic, fold, plateau, fault-block, and dome. A more detailed classification useful on a local scale predates plate tectonics and adds to these categories.
Volcanic mountains
Movements of tectonic plates create volcanoes along the plate boundaries, which erupt and form mountains. A volcanic arc system is a series of volcanoes that form near a subduction zone where the crust of a sinking oceanic plate melts and drags water down with the subducting crust.
Most volcanoes occur in a band encircling the Pacific Ocean (the Pacific Ring of Fire), and in another that extends from the Mediterranean across Asia to join the Pacific band in the Indonesian Archipelago. The most important types of volcanic mountain are composite cones or stratovolcanoes and shield volcanoes.
A shield volcano has a gently sloping cone because of the low viscosity of the emitted material, primarily basalt. Mauna Loa is the classic example, with a slope of 4°-6°. (The relation between slope and viscosity falls under the topic of angle of repose.) A composite volcano or stratovolcano has a more steeply rising cone (33°-40°), because of the higher viscosity of the emitted material, and eruptions are more violent and less frequent than for shield volcanoes. Examples include Vesuvius, Kilimanjaro, Mount Fuji, Mount Shasta, Mount Hood and Mount Rainier.
Fold mountains
When plates collide or undergo subduction (that is, ride one over another), the plates tend to buckle and fold, forming mountains. While volcanic arcs form at oceanic-continental plate boundaries, folding occurs at continental-continental plate boundaries. Most of the major continental mountain ranges are associated with thrusting and folding or orogenesis. Examples are the Balkan Mountains, the Jura and the Zagros mountains.
Block mountains
When a fault block is raised or tilted, a block mountain can result. Higher blocks are called horsts, and troughs are called grabens. A spreading apart of the surface causes tensional forces. When the tensional forces are strong enough to cause a plate to split apart, it does so such that a center block drops down relative to its flanking blocks.
An example is the Sierra Nevada range, where delamination created a block 650 km long and 80 km wide that consists of many individual portions tipped gently west, with east facing slips rising abruptly to produce the highest mountain front in the continental United States.
Another example is the Rila–Rhodope massif in Bulgaria, including the well defined horsts of Belasitsa (linear horst), Rila mountain (vaulted domed shaped horst) and Pirin mountain—a horst forming a massive anticline situated between the complex graben valleys of the Struma and Mesta rivers.
Uplifted passive margins
Unlike orogenic mountains there is no widely accepted geophysical model that explains elevated passive continental margins such as the Scandinavian Mountains, eastern Greenland, the Brazilian Highlands, or Australia's Great Dividing Range.
Different elevated passive continental margins most likely share the same mechanism of uplift. This mechanism is possibly related to far-field stresses in Earth's lithosphere. According to this view elevated passive margins can be likened to giant anticlinal lithospheric folds, where folding is caused by horizontal compression acting on a thin to thick crust transition zone (as are all passive margins).
Models
Hotspot volcanoes
Hotspots are supplied by a magma source in the Earth's mantle called a mantle plume. Although originally attributed to a melting of subducted oceanic crust, recent evidence belies this connection. The mechanism for plume formation remains a research topic.
Fault blocks
Several movements of the Earth's crust that lead to mountains are associated with faults. These movements actually are amenable to analysis that can predict, for example, the height of a raised block and the width of an intervening rift between blocks using the rheology of the layers and the forces of isostasy. Early bent plate models predicting fractures and fault movements have evolved into today's kinematic and flexural models.
| Physical sciences | Montane landforms | Earth science |
727455 | https://en.wikipedia.org/wiki/Chopper%20%28motorcycle%29 | Chopper (motorcycle) | A chopper is a type of custom motorcycle which emerged in the US state of California in the late 1950s. A chopper employs modified steering angles and lengthened forks for a stretched-out appearance. They can be built from an original motorcycle which is modified ("chopped") or built from scratch. Some of the characteristic features of choppers are long front ends with extended forks often coupled with an increased rake angle, hardtail frames (frames without rear suspension), very tall "ape hanger" or very short "drag" handlebars, lengthened or stretched frames, and larger than stock front wheel. To be considered a chopper a motorcycle frame must be cut and welded at some point. I.e. the name chopper. The "sissy bar", a set of tubes that connect the rear fender with the frame, and which are often extended several feet high, is a signature feature on many choppers.
Two famous examples of the chopper are customised Harley-Davidsons, the "Captain America" and "Billy Bike", seen in the 1969 film Easy Rider.
History
The Bob-Job Era, 1946–1959
Before there were choppers, there was the bobber, a motorcycle that had been "bobbed", or relieved of excess weight by removing parts. With the intent of making the bike lighter and faster, the fenders would often be removed, or at least to make it look better in the eyes of a rider seeking a more minimalist ride.
An early example of a bobber is the 1940 Indian Sport Scout "Bob-Job" which toured in the 1998 The Art of the Motorcycle exhibition. Indian Scouts and Chiefs of the time came with large, heavily valanced fenders, nearly reaching the center of the wheel on the 1941 Indian Series 441, while racing bikes had tiny fenders or none at all. The large bikes exemplified the "dresser" motorcycle aesthetic and provided a counterpoint to the minimalist bobber, and café racers.
In the post–World War II United States, servicemen returning home from the war started removing all parts deemed too big, heavy, ugly, or not essential to the basic function of the motorcycle, such as fenders, turn indicators, and even front brakes. The large, spring-suspended saddles were also removed in order to sit as low as possible on the motorcycle's frame. These machines were lightened to improve performance for dirt-track racing and mud racing. In California, dry lake beds were used for long top speed runs. Motorcycles and automobiles ran at the same meets, and bobbers were an important part of the hotrod culture that developed in this era.
The first choppers were built in America and were an outgrowth of the milder customization trend that had originated after WW2 when returning soldiers and others began modifying cars and motorcycles, frequently to improve performance in top-speed races on dry lake beds in Southern California and similar desolate spaces such as unused airstrips in other parts of the country, or on the street for street racing. These early modified motorcycles were known as "bobbers", and there are many common features between bobbers and choppers, with choppers differentiated by more radical modifications, especially frame tube and geometry modifications ("chopped" by welding) intended to make the bike longer.
The earliest choppers tended to be based on Harley-Davidson motorcycles, at first making use of the Flathead, Knucklehead and Panhead engines—many of which could be found in surplus military and police motorcycles bought cheaply at auction. As new engines became available, they were soon utilized in choppers. British bikes, particularly Triumphs, were also a popular motor for choppers early on. As the Japanese manufacturers began offering larger engines in the late 1960s these motors were also quickly put to use by chopper builders. The Honda 750-4 was the most widely used Japanese motor for early chopper builders. Choppers have been created using almost every available engine, but builders have always shown a preference for older air-cooled designs. It is rare to see a chopper with a radiator.
Over time the choppers became more and more about achieving a certain look, rather than being primarily performance-oriented modifications. The modifications that had had their origin in hot-rodding evolved into an artistic and aesthetic direction. By the mid-1970s stock Japanese and European performance motorcycles would outperform most bobbers and choppers. The one exception to this was the drag racing arena, which placed a premium on pure engine power, rather than handling over curvy courses. Chopper styling continued to be influenced by drag-bike modifications throughout the 1960s and 1970s.
While all choppers are highly customized bikes, sometimes even being built from scratch using all custom parts, not all customized bikes are choppers. In Europe at roughly the same era that choppers were invented and popularized in the US, bikers modified their bikes (primarily English brands like Triumph, BSA, Norton, and Matchless) in a different way, to achieve different looks, performance goals and riding position. The resulting bikes are known as café racers and look very different from a chopper.
As the popularity of choppers grew, in part through exposure in movies such as the 1969 classic Easy Rider, several motorcycle brands took note and began to include chopper influenced styling in their factory offerings. None of the factories were willing to go all out and do things like abandon rear-suspension to achieve the classic chopper look, however. As a result, these bikes were given the name "factory customs" and are not considered choppers.
Over the decades since the first choppers were created many different trends and fads have taken hold and held sway, so that it is often possible for someone to look at a chopper and say that it is a "1970s" style or fits into a specific era or sub-type. Some contemporary builders specialize in building choppers that very exactly fit into these styles, which are frequently referred to as "old school" style choppers.
Late 1950s to 1960s - early choppers
By the early 1960s there was a big enough contingent of people modifying motorcycles, still mostly big Harley-Davidsons, that a certain style had begun to take hold. A set of modifications became common: the fat tires and 16" wheels of the stock motorcycles were replaced with narrower tires often on a larger 19" or 21" wheel. Forward-mounted foot pegs replaced the standard large 'floorboard' footrests. Frequently the standard headlight and fuel tank were replaced with much smaller ones. Often upgraded chromed parts (either one-off fabricated replacements or manually chromed stock parts) were added. It is in this era that what we would today consider a chopper came into existence and began to be called the chopper.
During the 1960s, candy-colored paint, often multicolored and metal-flaked with different patterns, became a trend that allowed builders to further express their individuality and artistry. Soon many parts were being offered by small companies expressly for use in building choppers, not necessarily as performance parts as was common in the Bobber Era.
The first famous chopper builders came to prominence in this era, including Arlen Ness who was a leader in the "Frisco" or "Bay Area Chopper" style. Ness's bikes were characterized by having long low frames and highly raked front ends, typically 45 degrees or more, and frequently made use of springer front ends. Many made use of the newer Harley-Davidson Sportster motor, a simpler and more compact "unit motor" that included the transmission in the same housing as the motor itself, which lent itself nicely to Ness's stripped-down style. Many of Ness's bikes in this era retained the rear shocks of the donor Sportster to provide a more forgiving ride than the typical hardtail chopper.
In 1967 Denver Mullins (1944–1992) and Armando ″Mondo″ Porras opened "Denver's Choppers" in San Bernardino, California, and soon became famous for building "long bikes", often referred to as "Denver choppers". These featured even longer front ends than the Bay Area style and had a much higher frame (stretched "up and out"). Denver's was particularly well known for the springer forks that they fabricated, as well as the overall style of their bikes.
With choppers still not yet a mass market concept regional variations and formulations flourished in this period. Many innovations were tried in this period, found not to work that well, and then abandoned. A great deal of knowledge about how to build long bikes that handled well adjusting rake and trail was developed, yet less sophisticated builders also created a lot of bikes that had handling issues in this period as expertise was still scarce and closely held.
The 1970s: iconic choppers, diggers and Japanese motors
The huge success of the 1969 film Easy Rider instantly popularized the chopper around the world, and drastically increased the demand for them. What had been a subculture known to a relatively small group of enthusiasts in a few regions of the US became a global phenomenon. During the late 1960s, the first wave of European chopper builders emerged, such as the "Swedish Chopper" style, but Easy Rider brought attention everywhere to choppers.
The number of chopper-building custom shops multiplied, as did the number of suppliers of ready-made chopper parts. According to the taste and purse of the owner, chop shops would build high handlebars, or later Ed Roth's Wild Child designed stretched, narrowed, and raked front forks. Shops also custom-built exhaust pipes and many of the aftermarket kits followed in the late 1960s into the 1970s. Laws required (and in many locales still do) a retention fixture for the passenger, so vertical backrests called sissy bars became a popular installation, often sticking up higher than the rider's head.
While the decreased weight and lower seat position improved handling and performance, the main reason to build a chopper was to show off and provoke others by riding a machine that was stripped and almost nude compared to the stock Harley-Davidsons and automobiles of the period. Style trumped practicality, particularly as forks became longer and longer handling suffered. As one biker said, "You couldn't turn very good, but you sure looked good doing it."
The Digger became another popular style. Similar to the Frisco choppers Diggers were frequently even longer than earlier bikes, but still low. The coffin and prism shaped tanks on these bikes were frequently mated with very long front ends (12" over stock and more), with the archaic girder fork often being used to accomplish this instead of the more common springer or telescopic types. Body work was also moulded to flow seamlessly, using copious amounts of bondo. New paint colors and patterns included paisleys, day-glo and fluorescent, along with continuing use of metal-flakes and pearls.
Honda's groundbreaking 750 cc four-cylinder engine, first introduced to America in the 1969 CB-750, became widely available from salvage and wrecking operations and became a popular alternative to Harley-Davidson's motors. Harley's then-current big-twin motor, the Shovelhead, was extremely popular with chopper builders in this era, and use of the older motors, particularly the Knucklehead and Flathead declined as parts became harder to get and the performance of the new motors proved superior.
The 1980s and 1990s: improved engineering and aftermarket suppliers
In 1984 Harley-Davidson, who had been using chopper inspired styling for a number of years, released the 'Softail', a design that hid the rear shocks under the engine creating a profile that looked a lot like a hard tail. This frame was initially offered in the Softail Custom, a bike that took many styling cues from choppers, including the narrow 21" front wheel. Buyers looking for the chopper look had a plausible factory alternative, and interest in choppers declined.
With some time out of the limelight chopper builders seemed to work on craft more than wild innovation in this period. While individual builders still built long bikes, the trend was towards more moderate geometries, and the basics of how to build a good handling chopper while still looking great became more common knowledge. In this period, it became possible to assemble a complete chopper using all aftermarket parts, companies like S&S Cycle built complete replacement engines based on Harley-Davidson engines, frame makers such as Paughco offered a variety of hardtail frames, and many bikes were built using these new repro parts. Super long girder and springer forked bikes were less popular in this era, while the use of telescopic forks grew, and builders upgraded to larger diameter tubes in both forks and frames to gain more rigidity.
Japanese bike builders offered a dizzying array of new bikes, including full-faired racing styled machines as well as many 'customs' that picked at chopper styling in a random way and rarely achieved the powerful integrated style that more and more custom chopper builders in this era seemed able to consistently achieve. As materials, fabrication and knowledge improved the performance of the better choppers improved. More powerful engines drove the need for stronger frames, brakes and bigger tires with more grip. These trends worked together so that as the 1990s closed the modern chopper was a larger looking, more powerful machine. The widespread use of CNC made it possible for even small shops to fabricate out of block aluminum, and billet components became a signature item often replacing stamped and chromed steel components of the earlier eras.
The 21st Century: Reality television
The millennium began with the cable TV network The Discovery Channel creating a number of television shows around several custom bike builders who built choppers. The first, the 2000 special Motorcycle Mania, followed builder Jesse James of Long Beach, California, and is credited with creating "a new genre of reality TV" around choppers.
The celebrity builders featured on the cable shows enjoyed a large following. Companies like Jesse James' West Coast Choppers have been successful in producing expensive choppers, and a wide range of chopper-themed brands of merchandise such as clothing, automobile accessories and stickers.
The American Chopper reality television series featuring Paul Teutul Sr, and his sons Paul Jr. and Mike, ran six years starting in 2003, and featured bike building at Orange County Choppers (OCC).
2010: Backlash, Bobbers and the Old School Revival
This led to a backlash, and a renewed interest in home garage fabricated bikes built on a budget with available materials. Many builders eschewed Harley "pattern" motors and frames and started building choppers out of neglected bikes like Yamaha XS-650 twins, old Harley Sportsters, and various 1980's so called UJM bikes (four cylinder air-cooled Japanese bikes - Universal Japanese Motorcycle).
Another aspect of the backlash was a return to more traditional styling. Bobbers were again in style: stock rake machines with a stripped-down look, often with flat or primer paints in charcoal grey, flat black, olive drab or brown.
Indian Larry and Paul Cox along with other New York City builders from Psycho Cycles are often credited with leading the movement back towards old school choppers in this period. Indian Larry was a featured builder early on the series "Biker Build-Off" on Discovery network, and won all three build off competitions, highlighting the popularity of his old-school style.
Three-inch-wide belt drives and motors were still appreciated by many, but an increasing countermovement of people building bikes with Shovelhead motors and chain drive primaries has occurred. Springers and even girder forks have made yet another come back. Magazines such as Iron Horse, Street Chopper and Show Class cater to the retro, old-school and backyard builders, and feature more DIY technology than the TV builders with their million-dollar garages of the previous decade.
2020: Narrow Muscle Choppers
Later-generation builders take the chopper concept but keep the bike small, nimble and performance-inspired while nodding the styling of yesteryear. Some common characteristics of these custom bikes may include tall front end, narrow tires, high-output motor, cradle seat slammed onto the frame, t-bars, mid controls for example.
Choppers in the UK
In the UK, due to the cost and lack of availability of the v-twin engine, many chose to use British engines from bikes such as Triumph or BSA; following an increase in imports, Japanese engines have seen more use.
| Technology | Motorized road transport | null |
727811 | https://en.wikipedia.org/wiki/Snark%20%28graph%20theory%29 | Snark (graph theory) | In the mathematical field of graph theory, a snark is an undirected graph with exactly three edges per vertex whose edges cannot be colored with only three colors. In order to avoid trivial cases, snarks are often restricted to have additional requirements on their connectivity and on the length of their cycles. Infinitely many snarks exist.
One of the equivalent forms of the four color theorem is that every snark is a non-planar graph. Research on snarks originated in Peter G. Tait's work on the four color theorem in 1880, but their name is much newer, given to them by Martin Gardner in 1976. Beyond coloring, snarks also have connections to other hard problems in graph theory: writing in the Electronic Journal of Combinatorics, Miroslav Chladný and Martin Škoviera state that
As well as the problems they mention, W. T. Tutte's snark conjecture concerns the existence of Petersen graphs as graph minors of snarks; its proof has been long announced but remains unpublished, and would settle a special case of the existence of nowhere zero 4-flows.
History and examples
Snarks were so named by the American mathematician Martin Gardner in 1976, after the mysterious and elusive object of the poem The Hunting of the Snark by Lewis Carroll. However, the study of this class of graphs is significantly older than their name. Peter G. Tait initiated the study of snarks in 1880, when he proved that the four color theorem is equivalent to the statement that no snark is planar. The first graph known to be a snark was the Petersen graph; it was proved to be a snark by Julius Petersen in 1898, although it had already been studied for a different purpose by Alfred Kempe in 1886.
The next four known snarks were
the Blanuša snarks (two with 18 vertices), discovered by Danilo Blanuša in 1946,
the Descartes snark (210 vertices), discovered by Bill Tutte in 1948, and
the Szekeres snark (50 vertices), discovered by George Szekeres in 1973.
In 1975, Rufus Isaacs generalized Blanuša's method to construct two infinite families of snarks: the flower snarks and the Blanuša–Descartes–Szekeres snarks, a family that includes the two Blanuša snarks, the Descartes snark and the Szekeres snark. Isaacs also discovered a 30-vertex snark that does not belong to the Blanuša–Descartes–Szekeres family and that is not a flower snark: the double-star snark. Another infinite family, the Loupekine snarks, was published by Isaacs in 1976, credited to F. Loupekine. It includes two 22-vertex snarks derived from the Petersen graph.
The 50-vertex Watkins snark was discovered in 1989.
Another notable cubic non-three-edge-colorable graph is Tietze's graph, with 12 vertices; as Heinrich Franz Friedrich Tietze discovered in 1910, it forms the boundary of a subdivision of the Möbius strip requiring six colors. However, because it contains a triangle, it is not generally considered a snark. Under strict definitions of snarks, the smallest snarks are the Petersen graph and Blanuša snarks, followed by six different 20-vertex snarks.
A list of all of the snarks up to 36 vertices (according to a strict definition), and up to 34 vertices (under a weaker definition), was generated by Gunnar Brinkmann, Jan Goedgebeur, Jonas Hägglund and Klas Markström in 2012. The number of snarks for a given even number of vertices grows at least exponentially in the number of vertices. (Because they have odd-degree vertices, all snarks must have an even number of vertices by the handshaking lemma.) OEIS sequence contains the number of non-trivial snarks of vertices for small values of .
Definition
The precise definition of snarks varies among authors, but generally refers to cubic graphs (having exactly three edges at each vertex) whose edges cannot be colored with only three colors. By Vizing's theorem, the number of colors needed for the edges of a cubic graph is either three ("class one" graphs) or four ("class two" graphs), so snarks are cubic graphs of class two. However, in order to avoid cases where a snark is of class two for trivial reasons, or is constructed in a trivial way from smaller graphs, additional restrictions on connectivity and cycle lengths are often imposed. In particular:
If a cubic graph has a bridge, an edge whose removal would disconnect it, then it cannot be of class one. By the handshaking lemma, the subgraphs on either side of the bridge have an odd number of vertices each. Whichever of three colors is chosen for the bridge, their odd number of vertices prevents these subgraphs from being covered by cycles that alternate between the other two colors, as would be necessary in a 3-edge-coloring. For this reason, snarks are generally required to be bridgeless.
A loop (an edge connecting a vertex to itself) cannot be colored without causing the same color to appear twice at that vertex, a violation of the usual requirements for graph edge coloring. Additionally, a cycle consisting of two vertices connected by two edges can always be replaced by a single edge connecting their two other neighbors, simplifying the graph without changing its three-edge-colorability. For these reasons, snarks are generally restricted to simple graphs, graphs without loops or multiple adjacencies.
If a graph contains a triangle, then it can again be simplified without changing its three-edge-colorability, by contracting the three vertices of the triangle into a single vertex. Therefore, many definitions of snarks forbid triangles. However, although this requirement was also stated in Gardner's work giving the name "snark" to these graphs, Gardner lists Tietze's graph, which contains a triangle, as being a snark.
If a graph contains a four-vertex cycle, it can be simplified in two different ways by removing two opposite edges of the cycle and replacing the resulting paths of degree-two vertices by single edges. It has a three-edge-coloring if and only if at least one of these simplifications does. Therefore, Isaacs requires a "nontrivial" cubic class-two graph to avoid four-vertex cycles, and other authors have followed suit in forbidding these cycles. The requirement that a snark avoid cycles of length four or less can be summarized by stating that the girth of these graphs, the length of their shortest cycles, is at least five.
More strongly, the definition used by requires snarks to be cyclically 4-edge-connected. That means there can be no subset of three or fewer edges, the removal of which would disconnect the graph into two subgraphs each of which has at least one cycle. Brinkmann et al. define a snark to be a cubic and cyclically 4-edge-connected graph of girth five or more and class two; they define a "weak snark" to allow girth four.
Although these definitions only consider constraints on the girth up to five, snarks with arbitrarily large girth exist.
Properties
Work by Peter G. Tait established that the four-color theorem is true if and only if every snark is non-planar. This theorem states that every planar graph has a graph coloring of its the vertices with four colors, but Tait showed how to convert 4-vertex-colorings of maximal planar graphs into 3-edge-colorings of their dual graphs, which are cubic and planar, and vice versa. A planar snark would therefore necessarily be dual to a counterexample to the four-color theorem. Thus, the subsequent proof of the four-color theorem also demonstrates that all snarks are non-planar.
All snarks are non-Hamiltonian: when a cubic graph has a Hamiltonian cycle, it is always possible to 3-color its edges, by using two colors in alternation for the cycle, and the third color for the remaining edges. However, many known snarks are close to being Hamiltonian, in the sense that they are hypohamiltonian graphs: the removal of any single vertex leaves a Hamiltonian subgraph. A hypohamiltonian snark must be bicritical: the removal of any two vertices leaves a three-edge-colorable subgraph. The oddness of a cubic graph is defined as the minimum number of odd cycles, in any system of cycles that covers each vertex once (a 2-factor). For the same reason that they have no Hamiltonian cycles, snarks have positive oddness: a completely even 2-factor would lead to a 3-edge-coloring, and vice versa. It is possible to construct infinite families of snarks whose oddness grows linearly with their numbers of vertices.
The cycle double cover conjecture posits that in every bridgeless graph one can find a collection of cycles covering each edge twice, or equivalently that the graph can be embedded onto a surface in such a way that all faces of the embedding are simple cycles. When a cubic graph has a 3-edge-coloring, it has a cycle double cover consisting of the cycles formed by each pair of colors. Therefore, among cubic graphs, the snarks are the only possible counterexamples. More generally, snarks form the difficult case for this conjecture: if it is true for snarks, it is true for all graphs. In this connection, Branko Grünbaum conjectured that no snark could be embedded onto a surface in such a way that all faces are simple cycles and such that every two faces either are disjoint or share only a single edge; if any snark had such an embedding, its faces would form a cycle double cover. However, a counterexample to Grünbaum's conjecture was found by Martin Kochol.
Determining whether a given cyclically 5-connected cubic graph is 3-edge-colorable is NP-complete. Therefore, determining whether a graph is a snark is co-NP-complete.
Snark conjecture
W. T. Tutte conjectured that every snark has the Petersen graph as a minor. That is, he conjectured that the smallest snark, the Petersen graph, may be formed from any other snark by contracting some edges and deleting others. Equivalently (because the Petersen graph has maximum degree three) every snark has a subgraph that can be formed from the Petersen graph by subdividing some of its edges. This conjecture is a strengthened form of the four color theorem, because any graph containing the Petersen graph as a minor must be nonplanar. In 1999, Neil Robertson, Daniel P. Sanders, Paul Seymour, and Robin Thomas announced a proof of this conjecture. Steps towards this result have been published in 2016 and 2019, but the complete proof remains unpublished. See the Hadwiger conjecture for other problems and results relating graph coloring to graph minors.
Tutte also conjectured a generalization to arbitrary graphs: every bridgeless graph with no Petersen minor has a nowhere zero 4-flow. That is, the edges of the graph may be assigned a direction, and a number from the set {1, 2, 3}, such that the sum of the incoming numbers minus the sum of the outgoing numbers at each vertex is divisible by four. As Tutte showed, for cubic graphs such an assignment exists if and only if the edges can be colored by three colors, so the conjecture would follow from the snark conjecture in this case. However, proving the snark conjecture would not settle the question of the existence of 4-flows for non-cubic graphs.
| Mathematics | Graph theory | null |
728487 | https://en.wikipedia.org/wiki/Pie%20chart | Pie chart | A pie chart (or a circle chart) is a circular statistical graphic which is divided into slices to illustrate numerical proportion. In a pie chart, the arc length of each slice (and consequently its central angle and area) is proportional to the quantity it represents. While it is named for its resemblance to a pie which has been sliced, there are variations on the way it can be presented. The earliest known pie chart is generally credited to William Playfair's Statistical Breviary of 1801.
Pie charts are very widely used in the business world and the mass media. However, they have been criticized, and many experts recommend avoiding them, as research has shown it is difficult to compare different sections of a given pie chart, or to compare data across different pie charts. Pie charts can be replaced in most cases by other plots such as the bar chart, box plot, dot plot, etc.
History
The earliest known pie chart is generally credited to William Playfair's Statistical Breviary of 1801, in which two such graphs are used. Playfair presented an illustration, which contained a series of pie charts. One of those charts depicted the proportions of the Turkish Empire located in Asia, Europe and Africa before 1789. This invention was not widely used at first.
Playfair thought that pie charts were in need of a third dimension to add additional information.
Florence Nightingale may not have invented the pie chart, but she adapted it to make it more readable, which fostered its wide use, still today. Nightingale reconfigured the pie chart making the length of the wedges variable instead of their width. The graph, then, resembled a cock's comb. She was later assumed to have created it due to the obscurity and lack of practicality of Playfair's creation. Nightingale's polar area diagram, or occasionally the Nightingale rose diagram, equivalent to a modern circular histogram, to illustrate seasonal sources of patient mortality in the military field hospital she managed, was published in | Mathematics | Statistics | null |
729317 | https://en.wikipedia.org/wiki/Pathogen%20transmission | Pathogen transmission | In medicine, public health, and biology, transmission is the passing of a pathogen causing communicable disease from an infected host individual or group to a particular individual or group, regardless of whether the other individual was previously infected. The term strictly refers to the transmission of microorganisms directly from one individual to another by one or more of the following means:
airborne transmission – very small dry and wet particles that stay in the air for long periods of time allowing airborne contamination even after the departure of the host. Particle size < 5 μm.
droplet transmission – small and usually wet particles that stay in the air for a short period of time. Contamination usually occurs in the presence of the host. Particle size > 5 μm.
direct physical contact – touching an infected individual, including sexual contact
indirect physical contact – usually by touching a contaminated surface, including soil (fomite)
fecal–oral transmission – usually from unwashed hands, contaminated food or water sources due to lack of sanitation and hygiene, an important transmission route in pediatrics, veterinary medicine and developing countries.
via contaminated hypodermic needles or blood products
Transmission can also be indirect, via another organism, either a vector (e.g. a mosquito or fly) or an intermediate host (e.g. tapeworm in pigs can be transmitted to humans who ingest improperly cooked pork). Indirect transmission could involve zoonoses or, more typically, larger pathogens like macroparasites with more complex life cycles. Transmissions can be autochthonous (i.e. between two individuals in the same place) or may involve travel of the microorganism or the affected hosts.
A 2024 World Health Organization report standardized the terminology for the transmission modes of all respiratory pathogens in alignment with particle physics: airborne transmission; inhalation; direct deposition; and contact. But these newly standardized terms have yet to be translated to policy, including infection control policy or the pandemic accords or updated International Health Regulations.
Definition and related terms
An infectious disease agent can be transmitted in two ways: as disease agent transmission from one individual to another in the same generation (peers in the same age group) by either direct contact (licking, touching, biting), or indirect contact through air – cough or sneeze (vectors or fomites that allow the transmission of the agent causing the disease without physical contact)
or by vertical disease transmission, passing the agent causing the disease from parent to offspring, such as in prenatal or perinatal transmission.
The term infectivity describes the ability of an organism to enter, survive and multiply in the host, while the infectiousness of a disease agent indicates the comparative ease with which the disease agent is transmitted to other hosts. Transmission of pathogens can occur by direct contact, through contaminated food, body fluids or objects, by airborne inhalation or through vector organisms.
Transmissibility is the probability of an infection, given a contact between an infected host and a noninfected host.
Community transmission means that the source of infection for the spread of an illness is unknown or a link in terms of contacts between patients and other people is missing. It refers to the difficulty in grasping the epidemiological link in the community beyond confirmed cases.
Local transmission means that the source of the infection has been identified within the reporting location (such as within a country, region or city).
Routes of transmission
The route of transmission is important to epidemiologists because patterns of contact vary between different populations and different groups of populations depending on socio-economic, cultural and other features. For example, low personal and food hygiene due to the lack of a clean water supply may result in increased transmission of diseases by the fecal-oral route, such as cholera. Differences in incidence of such diseases between different groups can also throw light on the routes of transmission of the disease. For example, if it is noted that polio is more common in cities in underdeveloped countries, without a clean water supply, than in cities with a good plumbing system, we might advance the theory that polio is spread by the fecal-oral route. Two routes are considered to be airborne: Airborne infections and droplet infections.
Airborne infection
"Airborne transmission refers to infectious agents that are spread via droplet nuclei (residue from evaporated droplets) containing infective microorganisms. These organisms can survive outside the body and remain suspended in the air for long periods of time. They infect others via the upper and lower respiratory tracts." The size of the particles for airborne infections need to be < 5 μm. It includes both dry and wet aerosols and thus requires usually higher levels of isolation since it can stay suspended in the air for longer periods of time. i.e., separate ventilation systems or negative pressure environments are needed to avoid general contamination. e.g., tuberculosis, chickenpox, measles.
infection
A common form of transmission is by way of respiratory droplets, generated by coughing, sneezing, or talking. Respiratory droplet transmission is the usual route for respiratory infections. Transmission can occur when respiratory droplets reach susceptible mucosal surfaces, such as in the eyes, nose or mouth. This can also happen indirectly via contact with contaminated surfaces when hands then touch the face. Before drying, respiratory droplets are large and cannot remain suspended in the air for long, and are usually dispersed over short distances. The size of the particles for droplet infections are > 5 μm.
Organisms spread by droplet transmission include respiratory viruses such as influenza virus, parainfluenza virus, adenoviruses, rhinovirus, respiratory syncytial virus, human metapneumovirus, Bordetella pertussis, pneumococci, streptococcus pyogenes, diphtheria, rubella, and coronaviruses. Spread of respiratory droplets from the wearer can be reduced through wearing of a surgical mask.
Direct contact
Direct contact occurs through skin-to-skin contact, kissing, and sexual intercourse. Direct contact also refers to contact with soil or vegetation harboring infectious organisms. Additionally, while fecal–oral transmission is primarily considered an indirect contact route, direct contact can also result in transmission through feces.
Diseases that can be transmitted by direct contact are called contagious (contagious is not the same as infectious; although all contagious diseases are infectious, not all infectious diseases are contagious). These diseases can also be transmitted by sharing a towel (where the towel is rubbed vigorously on both bodies) or items of clothing in close contact with the body (socks, for example) if they are not washed thoroughly between uses. For this reason, contagious diseases often break out in schools, where towels are shared and personal items of clothing accidentally swapped in the changing rooms.
Some diseases that are transmissible by direct contact include athlete's foot, impetigo, syphilis, warts, and conjunctivitis.
Sexual
This refers to any infection that can be caught during sexual activity with another person, including vaginal or anal sex, less commonly through oral sex (see below) and rarely through manual sex (see below). Transmission is either directly between surfaces in contact during intercourse (the usual route for bacterial infections and those infections causing sores) or from secretions (semen or the fluid secreted by the excited female) which carry infectious agents that get into the partner's blood stream through tiny tears in the penis, vagina or rectum (this is a more usual route for viruses). In this second case, anal sex is considerably more hazardous since the penis opens more tears in the rectum than the vagina, as the vagina is more elastic and more accommodating.
Some infections transmissible by the sexual route include HIV/AIDS, chlamydia, genital warts, gonorrhea, hepatitis B, syphilis, herpes, and trichomoniasis.
Oral sex
Sexually transmitted infections such as HIV and hepatitis B are thought to not normally be transmitted through mouth-to-mouth contact, although it is possible to transmit some STIs between the genitals and the mouth, during oral sex. In the case of HIV, this possibility has been established. It is also responsible for the increased incidence of herpes simplex virus 1 (which is usually responsible for oral infections) in genital infections and the increased incidence of the type 2 virus (more common genitally) in oral infections.
Manual sex
While rare in regards to this sexual practice, some infections that can spread via manual sex include HPV, chlamydia, and syphilis.
Oral
Infections that are transmitted primarily by oral means may be caught through direct oral contact such as kissing, or by indirect contact such as by sharing a drinking glass or a cigarette. Infections that are known to be transmissible by kissing or by other direct or indirect oral contact include all of the infections transmissible by droplet contact and (at least) all forms of herpes viruses, namely Cytomegalovirus infections herpes simplex virus (especially HSV-1) and infectious mononucleosis.
Mother-to-child transmission
This is from mother to child (more rarely father to child), often in utero, during childbirth (also referred to as perinatal infection) or during postnatal physical contact between parents and offspring. In mammals, including humans, it occurs also via breast milk (transmammary transmission). Infectious diseases that can be transmitted in this way include: HIV, hepatitis B and syphilis. Many mutualistic organisms are transmitted vertically.
Iatrogenic
Transmission due to medical procedures, such as touching a wound, the use of contaminated medical equipment, or an injection or transplantation of infected material. Some diseases that can be transmitted iatrogenically include Creutzfeldt–Jakob disease, HIV, and many more.
Needle sharing
This is the practice of intravenous drug-users by which a needle or syringe is shared by multiple individuals to administer intravenous drugs such as heroin, steroids, and hormones. This can act as a vector for blood-borne diseases, such as Hepatitis C (HCV) and HIV.
Indirect contact
Indirect contact transmission, also known as vehicle-borne transmission, involves transmission through contamination of inanimate objects. Vehicles that may indirectly transmit an infectious agent include food, water, biologic products such as blood, and fomites such as handkerchiefs, bedding, or surgical scalpels. A vehicle may passively carry a pathogen, as in the case of food or water may carrying hepatitis A virus. Alternatively, the vehicle may provide an environment in which the agent grows, multiplies, or produces toxin, such as improperly canned foods provide an environment that supports production of botulinum toxin by Clostridium botulinum.
Transmission by other organisms
A vector is an organism that does not cause disease itself but that transmits infection by conveying pathogens from one host to another.
Vectors may be mechanical or biological. A mechanical vector picks up an infectious agent on the outside of its body and transmits it in a passive manner. An example of a mechanical vector is a housefly, which lands on cow dung, contaminating its appendages with bacteria from the feces, and then lands on food prior to consumption. The pathogen never enters the body of the fly. In contrast, biological vectors harbor pathogens within their bodies and deliver pathogens to new hosts in an active manner, usually a bite. Biological vectors are often responsible for serious blood-borne diseases, such as malaria, viral encephalitis, Chagas disease, Lyme disease and African sleeping sickness. Biological vectors are usually, though not exclusively, arthropods, such as mosquitoes, ticks, fleas and lice. Vectors are often required in the life cycle of a pathogen. A common strategy used to control vector-borne infectious diseases is to interrupt the life cycle of a pathogen by killing the vector.
Fecal–oral
In the fecal-oral route, pathogens in fecal particles pass from one person to the mouth of another person. Although it is usually discussed as a route of transmission, it is actually a specification of the entry and exit portals of the pathogen, and can operate across several of the other routes of transmission. Fecal–oral transmission is primarily considered as an indirect contact route through contaminated food or water. However, it can also operate through direct contact with feces or contaminated body parts, such as through anal sex. It can also operate through droplet or airborne transmission through the toilet plume from contaminated toilets.
Main causes of fecal–oral disease transmission include lack of adequate sanitation and poor hygiene practices - which can take various forms. Fecal oral transmission can be via foodstuffs or water that has become contaminated. This can happen when people do not adequately wash their hands after using the toilet and before preparing food or tending to patients.
The fecal-oral route of transmission can be a public health risk for people in developing countries who live in urban slums without access to adequate sanitation. Here, excreta or untreated sewage can pollute drinking water sources (groundwater or surface water). The people who drink the polluted water can become infected. Another problem in some developing countries, is open defecation which leads to disease transmission via the fecal-oral route.
Even in developed countries there are periodic system failures resulting in a sanitary sewer overflow. This is the typical mode of transmission for infectious agents such as cholera, hepatitis A, polio, Rotavirus, Salmonella, and parasites (e.g. Ascaris lumbricoides).
Tracking
Tracking the transmission of infectious diseases is called disease surveillance. Surveillance of infectious diseases in the public realm traditionally has been the responsibility of public health agencies, on an international, national, or local level. Public health staff relies on health care workers and microbiology laboratories to report cases of reportable diseases to them. The analysis of aggregate data can show the spread of a disease and is at the core of the specialty of epidemiology.
To understand the spread of the vast majority of non-notifiable diseases, data either need to be collected in a particular study, or existing data collections can be mined, such as insurance company data or antimicrobial drug sales for example.
For diseases transmitted within an institution, such as a hospital, prison, nursing home, boarding school, orphanage, refugee camp, etc., infection control specialists are employed, who will review medical records to analyze transmission as part of a hospital epidemiology program, for example.
Because these traditional methods are slow, time-consuming, and labor-intensive, proxies of transmission have been sought. One proxy in the case of influenza is tracking of influenza-like illness at certain sentinel sites of health care practitioners within a state, for example. Tools have been developed to help track influenza epidemics by finding patterns in certain web search query activity. It was found that the frequency of influenza-related web searches as a whole rises as the number of people sick with influenza rises. Examining space-time relationships of web queries has been shown to approximate the spread of influenza and dengue.
Computer simulations of infectious disease spread have been used.
Human aggregation can drive transmission, seasonal variation and outbreaks of infectious diseases, such as the annual start of school, bootcamp, the annual Hajj etc. Most recently, data from cell phones have been shown to be able to capture population movements well enough to predict the transmission of certain infectious diseases, like rubella.
Relationship with virulence and survival
Pathogens must have a way to be transmitted from one host to another to ensure their species' survival. Infectious agents are generally specialized for a particular method of transmission. Taking an example from the respiratory route, from an evolutionary perspective viruses or bacteria that cause their host to develop coughing and sneezing symptoms have a great survival advantage, as they are much more likely to be ejected from one host and carried to another. This is also the reason that many microorganisms cause diarrhea.
The relationship between virulence and transmission is complex and has important consequences for the long term evolution of a pathogen. Since it takes many generations for a microbe and a new host species to co-evolve, an emerging pathogen may hit its earliest victims especially hard. It is usually in the first wave of a new disease that death rates are highest. If a disease is rapidly fatal, the host may die before the microbe can be passed along to another host. However, this cost may be overwhelmed by the short-term benefit of higher infectiousness if transmission is linked to virulence, as it is for instance in the case of cholera (the explosive diarrhea aids the bacterium in finding new hosts) or many respiratory infections (sneezing and coughing create infectious aerosols).
Anything that reduces the rate of transmission of an infection carries positive externalities, which are benefits to society that are not reflected in a price to a consumer. This is recognized implicitly when vaccines are offered for free or at a cost to the patient less than the purchase price.
Beneficial microorganisms
The mode of transmission is also an important aspect of the biology of beneficial microbial symbionts, such as coral-associated dinoflagellates or human microbiota. Organisms can form symbioses with microbes transmitted from their parents, from the environment or unrelated individuals, or both.
Vertical transmission
Vertical transmission refers to acquisition of symbionts from parents (usually mothers). Vertical transmission can be intracellular (e.g. transovarial), or extracellular (for example through post-embryonic contact between parents and offspring). Both intracellular and extracellular vertical transmission can be considered a form of non-genetic inheritance or parental effect. It has been argued that most organisms experience some form of vertical transmission of symbionts. Canonical examples of vertically transmitted symbionts include the nutritional symbiont Buchnera in aphids (transovarially transmitted intracellular symbiont) and some components of the human microbiota (transmitted during passage of infants through the birth canal and also through breastfeeding).
Horizontal transmission
Some beneficial symbionts are acquired horizontally, from the environment or unrelated individuals. This requires that host and symbiont have some method of recognizing each other or each other's products or services. Often, horizontally acquired symbionts are relevant to secondary rather than primary metabolism, for example for use in defense against pathogens, but some primary nutritional symbionts are also horizontally (environmentally) acquired. Additional examples of horizontally transmitted beneficial symbionts include bioluminescent bacteria associated with bobtail squid and nitrogen-fixing bacteria in plants.
Mixed-mode transmission
Many microbial symbionts, including human microbiota, can be transmitted both vertically and horizontally. Mixed-mode transmission can allow symbionts to have the "best of both worlds" – they can vertically infect host offspring when host density is low, and horizontally infect diverse additional hosts when a number of additional hosts are available. Mixed-mode transmission make the outcome (degree of harm or benefit) of the relationship more difficult to predict, because the evolutionary success of the symbiont is sometimes but not always tied to the success of the host.
| Biology and health sciences | Concepts | Health |
729500 | https://en.wikipedia.org/wiki/Head%20and%20neck%20cancer | Head and neck cancer | Head and neck cancer is a general term encompassing multiple cancers that can develop in the head and neck region. These include cancers of the mouth, tongue, gums and lips (oral cancer), voice box (laryngeal), throat (nasopharyngeal, oropharyngeal, hypopharyngeal), salivary glands, nose and sinuses.
Head and neck cancer can present a wide range of symptoms depending on where the cancer developed. These can include an ulcer in the mouth that does not heal, changes in the voice, difficulty swallowing, red or white patches in the mouth, and a neck lump.
The majority of head and neck cancer is caused by the use of alcohol or tobacco (including smokeless tobacco). An increasing number of cases are caused by the human papillomavirus (HPV). Other risk factors include the Epstein–Barr virus, chewing betel quid (paan), radiation exposure, poor nutrition and workplace exposure to certain toxic substances. About 90% are pathologically classified as squamous cell cancers. The diagnosis is confirmed by a tissue biopsy. The degree of surrounding tissue invasion and distant spread may be determined by medical imaging and blood tests.
Not using tobacco or alcohol can reduce the risk of head and neck cancer. Regular dental examinations may help to identify signs before the cancer develops. The HPV vaccine helps to prevent HPV-related oropharyngeal cancer. Treatment may include a combination of surgery, radiation therapy, chemotherapy, and targeted therapy. In the early stage head and neck cancers are often curable but 50% of people see their doctor when they already have an advanced disease.
Globally, head and neck cancer accounts for 650,000 new cases of cancer and 330,000 deaths annually on average. In 2018, it was the seventh most common cancer worldwide, with 890,000 new cases documented and 450,000 people dying from the disease. The usual age at diagnosis is between 55 and 65 years old. The average 5-year survival following diagnosis in the developed world is 42–64%.
Signs and symptoms
Head and neck cancers can cause a broad range of symptoms, many of which occur together. These can be categorised local (head and neck cancer-specific), general and gastrointestinal symptoms. Local symptoms include changes in taste and voice, inflammation of the mouth or throat (mucositis), dry mouth (xerostomia), and difficulty swallowing (dysphagia). General symptoms include difficulty sleeping, tiredness, depression, nerve damage (peripheral neuropathy). Gastrointestinal symptoms are typically nausea and vomiting.
Symptoms predominantly include a sore on the face or oral cavity that does not heal, trouble swallowing, or a change in voice. In those with advanced disease, there may be unusual bleeding, facial pain, numbness or swelling, and visible lumps on the outside of the neck or oral cavity. Head and neck cancer often begins with benign signs and symptoms of the disease, like an enlarged lymph node on the outside of the neck, a hoarse-sounding voice, or a progressive worsening cough or sore throat. In the case of head and neck cancer, these symptoms will be notably persistent and become chronic. There may be a lump or a sore in the throat or neck that does not heal or go away. There may be difficulty or pain in swallowing. Speaking may become difficult. There may also be a persistent earache.
Other symptoms can include: a lump in the lip, mouth, or gums; ulcers or mouth sores that do not heal; bleeding from the mouth or numbness; bad breath; discolored patches that persist in the mouth; a sore tongue; and slurring of speech if the cancer is affecting the tongue. There may also be congested sinuses, weight loss, and some numbness or paralysis of facial muscles.
Mouth
Oral cancer affects the areas of the mouth, including the inner lip, tongue, floor of the mouth, gums, and hard palate. Cancers of the mouth are strongly associated with tobacco use, especially the use of chewing tobacco or dipping tobacco, as well as heavy alcohol use. Cancers of this region, particularly the tongue, are more frequently treated with surgery than other head and neck cancers. Lip and oral cavity cancers are the most commonly encountered types of head and neck cancer.
Surgeries for oral cancers include:
Maxillectomy (can be done with or without orbital exenteration)
Mandibulectomy (removal of the lower jaw or part of it)
Glossectomy (tongue removal; can be total, hemi, or partial)
Radical neck dissection
Combinational (e.g., glossectomy and laryngectomy done together).
The defect is typically covered or improved by using another part of the body and/or skin grafts and/or wearing a prosthesis.
Nose
Paranasal sinus and nasal cavity cancer affects the nasal cavity and the paranasal sinuses. Most of these cancers are squamous cell carcinomas.
Nasopharynx
Nasopharyngeal cancer arises in the nasopharynx, the region in which the nasal cavities and the Eustachian tubes connect with the upper part of the throat. While some nasopharyngeal cancers are biologically similar to the common head and neck cancers, "poorly differentiated" nasopharyngeal carcinoma is lymphoepithelioma, which is distinct in its epidemiology, biology, clinical behavior, and treatment and is treated as a separate disease by many experts.
Throat
Most oropharyngeal cancers begin in the oropharynx (throat), the middle part of the throat that includes the soft palate, the base of the tongue, and the tonsils. Cancers of the tonsils are more strongly associated with human papillomavirus infection than are cancers of other regions of the head and neck. HPV-positive oropharyngeal cancer generally has a better outcome than HPV-negative disease, with a 54% better survival rate, but this advantage for HPV-associated cancer applies only to oropharyngeal cancers.
People with oropharyngeal carcinomas are at high risk of developing a second primary head and neck cancer.
Hypopharynx
The hypopharynx includes the pyriform sinuses, the posterior pharyngeal wall, and the postcricoid area. Tumors of the hypopharynx frequently have an advanced stage at diagnosis and have the most adverse prognoses of pharyngeal tumors. They tend to metastasize early due to the extensive lymphatic network around the larynx.
Larynx
Laryngeal cancer begins in the larynx, or "voice box", and is the second most common type of head and neck cancer encountered. Cancer may occur on the vocal folds themselves ("glottic" cancer) or on tissues above and below the true cords ("supraglottic" and "subglottic" cancers, respectively). Laryngeal cancer is strongly associated with tobacco smoking.
Surgery can include laser excision of small vocal cord lesions, partial laryngectomy (removal of part of the larynx), or total laryngectomy (removal of the whole larynx). If the whole larynx has been removed, the person is left with a permanent tracheostomy. Voice rehabilitation in such patients can be achieved in three important ways: esophageal speech, tracheoesophageal puncture, or electrolarynx. One would likely require intensive teaching, speech therapy, and/or an electronic device.
Trachea and salivary glands
Cancer of the trachea is a rare cancer usually classified as a lung cancer.
Most tumors of the salivary glands differ from the common head and neck cancers in cause, histopathology, clinical presentation, and therapy. Other uncommon tumors arising in the head and neck include teratomas, adenocarcinomas, adenoid cystic carcinomas, and mucoepidermoid carcinomas. Rarer still are melanomas and lymphomas of the upper aerodigestive tract.
Causes
Alcohol and tobacco
Alcohol and tobacco use are major risk factors for head and neck cancer. 72% of head and neck cancer cases are caused by using both alcohol and tobacco. This rises to 89% when looking specifically at laryngeal cancer.
There is thought to be a dose-dependent relationship between alcohol use and development of head and neck cancer where higher rates of alcohol consumption contribute to an increased risk of developing head and neck cancer. Alcohol use following a diagnosis of head and neck cancer also contributes to other negative outcomes. These include physical effects such as an increased risk of developing a second primary cancer or other malignancies, cancer recurrence, and worse prognosis in addition to an increased chance of having a future feeding tube placed and osteoradionecrosis of the jaw. Negative social factors are also increased with sustained alcohol use after diagnosis including unemployment and work disability.
The way in which alcohol contributes to cancer development is not fully understood. It is thought to be related to permanent damage of DNA strands by a metabolite of alcohol called acetaldehyde. Other suggested mechanisms include nutritional deficiencies and genetic variations.
Tobacco smoking is one of the main risk factors for head and neck cancer. Cigarette smokers have a lifetime increased risk for head and neck cancer that is 5 to 25 times higher than the general population. The ex-smoker's risk of developing head and neck cancer begins to approach the risk in the general population 15 years after smoking cessation. In addition, people who smoke have a worse prognosis than those who have never smoked. Furthermore, people who continue to smoke after diagnosis of head and neck cancer have the highest probability of dying compared to those who have never smoked. This effect is seen in patients with HPV-positive head and neck cancer as well. It has also been demonstrated that passive smoking, both at work and at home, increases the risk of head and neck cancer.
A major carcinogenic compound in tobacco smoke is acrylonitrile. Acrylonitrile appears to indirectly cause DNA damage by increasing oxidative stress, leading to increased levels of 8-oxo-2'-deoxyguanosine (8-oxo-dG) and formamidopyrimidine in DNA. (see image). Both 8-oxo-dG and formamidopyrimidine are mutagenic. DNA glycosylase NEIL1 prevents mutagenesis by 8-oxo-dG and removes formamidopyrimidines from DNA.
Smokeless tobacco (including products where tobacco is chewed) is a cause of oral cancer. Increased risk of oral cancer caused by smokeless tobacco is present in countries such as the United States but particularly prevalent in Southeast Asian countries where the use of smokeless tobacco is common. Smokeless tobacco is associated with a higher risk of developing head and neck cancer due to the presence of the tobacco-specific carcinogen N'-nitrosonornicotine.
Cigar and pipe smoking are also important risk factors for oral cancer. They have a dose dependent relationship with more consumption leading to higher chances of developing cancer. The use of electronic cigarettes may also lead to the development of head and neck cancers due to the substances like propylene glycol, glycerol, nitrosamines, and metals contained therein, which can cause damage to the airways. Exposure to e-vapour has been shown to reduce cell viability and increase the rate of cell death via apoptosis or necrosis with or without nicotine. This area of study requires more research, however. Similarly, additional research is needed to understand how marijuana possibly promotes head and neck cancers. A 2019 meta-analysis did not conclude that marijuana was associated with head and neck cancer risk. Yet individuals with cannabis use disorder were more likely to be diagnosed with such cancers in a large study published 2024.
Diet
Many dietary nutrients are associated with cancer protection and its development. Generally, foods with a protective effect with respect to oral cancer demonstrate antioxidant and anti-inflammatory effects such as fruits, vegetables, curcumin and green tea. Conversely, pro-inflammatory food substances such as red meat, processed meat and fried food can increase the risk of developing head and neck cancer. An increased adherence to the Mediterranean diet is also related to a lower risk of cancer mortality and a reduced risk of developing multiple cancers including head and neck cancer. Elevated levels of nitrites in preserved meats and salted fish have been shown to increase the risk of nasopharyngeal cancer. Overall, a poor nutritional intake (often associated with alcoholism) with subsequent vitamin deficiencies is a risk factor for head and neck cancer.
In terms of nutritional supplements, antioxidants such as vitamin E and beta-carotene might reduce the toxic effect of radiotherapy in people with head and neck cancer but they can also increase recurrence rates, especially in smokers.
Betel nut
Betel nut chewing is associated with an increased risk of head and neck cancer. When chewed with additional tobacco in its preparation (like in gutka), there is an even higher risk, especially for oral and oropharyngeal cancers.
Genetics
People who develop head and neck cancer may have a genetic predisposition for the condition. There are seven known genetic variations (loci) which specifically increase the chances of developing oral and pharyngeal cancer. Family history, that is having a first-degree relative with head and neck cancer, is also a risk factor. In addition, genetic variations in pathways involved in alcohol metabolism (for example alcohol dehydrogenase) have been associated with an increased head and neck cancer risk.
Radiation
It is known that prior exposure to radiation of the head and neck is associated with an increased risk of cancer, particularly thyroid, salivary gland and squamous cell carcinomas, although there is a time-delay of many years and the overall risk is still low.
Infection
Human papillomavirus
Some head and neck cancers, and in particular oropharyngeal cancer, are caused by the human papillomavirus (HPV), and 70% of all head and neck cancer cases are related to HPV. Risk factors for HPV-positive oropharyngeal cancer include multiple sexual partners, anal and oral sex and a weak immune system.
The incidence of HPV-related head and neck cancer is increasing, especially in the Western world. In the United States, the incidence of HPV-positive oropharyngeal cancer has overtaken HPV-positive cervical cancer as the leading HPV related cancer type. An increased incidence has particularly affected males. As a result, recent changes have resulted in the HPV vaccine being offered to adolescent boys between 12-13 (previously only offered to girls between this age due to cervical cancer risks) and men under 45 who have sex with men in the UK.
Over 20 different high-risk HPV subtypes have been implicated in causing head and neck cancer. In particular, HPV-16 is responsible for up to 90% of oropharyngeal cancer in North America. Approximately 15–25% of head and neck cancers contain genomic DNA from HPV, and the association varies based on the site of the tumor. In the case of HPV-positive oropharyngeal cancer, the highest distribution is in the tonsils, where HPV DNA is found in 45–67% of the cases, and it is less often in the hypopharynx (13–25%), and least often in the oral cavity (12–18%) and larynx (3–7%).
Positive HPV16 status is associated with a improved prognosis over HPV-negative oropharyngeal cancer due to better response to radiotherapy and chemotherapy.
HPV can induce tumors by several mechanisms:
E6 and E7 oncogenic proteins.
Disruption of tumor suppressor genes.
High-level DNA amplifications, for example, oncogenes.
Generating alternative nonfunctional transcripts.
Interchromosomal rearrangements.
Distinct host genome methylation and expression patterns, produced even when the virus is not integrated into the host genome.
There are observed biological differences between HPV-positive and HPV-negative head and neck cancer, for example in terms of mutation patterns. In HPV-negative disease, genes frequently mutated include TP53, CDKN2A and PIK3CA. In HPV-positive disease, these genes are less frequently mutated, and the tumour suppressor gene p53 and pRb (protein retinoblastoma) are commonly inactivated by HPV oncoproteins E6 and E7 respectively. In addition, viral infections such as HPV can cause aberrant DNA methylation during cancer development. HPV-positive head and neck cancers demonstrate higher levels of such DNA methylation compared to HPV-negative disease.
E6 sequesters p53 to promote p53 degradation, while E7 inhibits pRb. Degradation of p53 results in cells being unable to respond to checkpoint signals that are normally present to activate apoptosis when DNA damage is signalled. Loss of pRb leads to deregulation of cell proliferation and apoptosis. Both mechanisms therefore leave cell proliferation unchecked and increase the chance of carcinogenesis.
Epstein–Barr virus
Epstein–Barr virus (EBV) infection is associated with nasopharyngeal cancer. Nasopharyngeal cancer caused by EBV commonly occurs in some countries of the Mediterranean and Asia, where EBV antibody titers can be measured to screen high-risk populations.
Gastroesophageal reflux disease
The presence of gastroesophageal reflux disease (GERD) or laryngeal reflux disease can also be a major factor. Stomach acids that flow up through the esophagus can damage its lining and raise susceptibility to throat cancer.
Hematopoietic stem cell transplantation
People after hematopoietic stem cell transplantation (HSCT) are at a higher risk for oral cancer. Post-HSCT oral cancer may have more aggressive behavior and a poorer prognosis when compared to oral cancer in non-HSCT patients. This effect is supposed to be due to continuous, lifelong immune suppression and chronic oral graft-versus-host disease.
Other risk factors
Several other risk factors have been identified in the development of head and neck cancer. These include occupational environmental carcinogen exposure such as asbestos, wood dust, mineral acid, sulfuric acid mists and metal dusts. In addition, weakened immune systems, age greater than 55 years, poor socioeconomic factors such as lower incomes and occupational status, and low body mass index (<18.5 kg/m2) are also risk factors. Poor oral hygiene and chronic oral cavity inflammation (for example secondary to chronic gum inflammation) are also linked to an increased head and neck cancer risk. The presence of leukoplakia, which is the appearance of white patches or spots in the mouth, can develop into cancer in about 1⁄3 of cases.
Diagnosis
A significant proportion of people with head and neck cancer will present to their physicians with an already advanced stage disease. This can either be down to patient factors (delays in seeking medical attention), or physician factors (such as delays in referral from primary care, or non-diagnostic investigation results).
A person usually presents to the physician complaining of one or more of the typical symptoms. These symptoms may be site specific (such as a laryngeal cancer causing hoarse voice), or not site specific (earache can be caused by multiple types of head and neck cancers).
The physician will undertake a thorough history to determine the nature of the symptoms and the presence or absence of any risk factors. The physician will also ask about other illnesses such as heart or lung diseases as they may impact their fitness for potentially curative treatment. Clinical examination will involve examination of the neck for any masses, examining inside the mouth for any abnormalities and assessing the rest of the pharynx and larynx with a nasendoscope.
Further investigations will be directed by the symptoms discussed and any abnormalities identified during the exam.
Neck masses typically undergo assessment with ultrasound and a fine-needle aspiration (FNA, a type of needle biopsy). Concerning lesions that are readily accessible (such as in the mouth) can be biopsied with a local anaesthetic. Lesions less readily available can be biopsied either with the patient awake or under a general anaesthetic depending on local expertise and availability of specialist equipment.
The cancer will also need to be staged (accurately determine its size, association with nearby structures, and spread to distant sites). This is typically done by scanning the patient with a combination of magnetic resonance imaging (MRI), computed tomography (CT) and/or positron emission tomography (PET). Exactly which investigations are required will depend on a variety of factors such as the site of concern and the size of the tumour.
Some people will present with a neck lump containing cancer cells (identified by FNA) that have spread from elsewhere, but with no identifiable primary site on initial assessment. In such cases people will undergo additional testing to attempt to find the initial site of cancer, as this has significant implications for their treatment. These patients undergo MRI scanning, PET-CT and then panendoscopy and biopsies of any abnormal areas. If the scans and panendoscopy still do not identify a primary site for the cancer, affected people will undergo a bilateral tonsillectomy and tongue base mucosectomy (as these are the most common subsites of cancer that spread to the neck). This procedure can be done with or without robotic assistance.
Once a diagnosis is confirmed, a multidisciplinary discussion of the optimal treatment strategy will be undertaken between the radiation oncologist, surgical oncologist, and medical oncologist. A histopathologist and a radiologist will also be present to discuss the biopsy and imaging findings. Most (90%) cancers of the head and neck are squamous cell-derived, termed "head-and-neck squamous-cell carcinomas".
Histopathology
Throat cancers are classified according to their histology or cell structure and are commonly referred to by their location in the oral cavity and neck. This is because where the cancer appears in the throat affects the prognosis; some throat cancers are more aggressive than others, depending on their location. The stage at which the cancer is diagnosed is also a critical factor in the prognosis of throat cancer. Treatment guidelines recommend routine testing for the presence of HPV for all oropharyngeal squamous cell carcinoma tumors.
Squamous-cell carcinoma
Squamous-cell carcinoma is a cancer of the squamous cell, a kind of epithelial cell found in both the skin and mucous membranes. It accounts for over 90% of all head and neck cancers, including more than 90% of throat cancer. Squamous cell carcinoma is most likely to appear in males over 40 years of age with a history of heavy alcohol use coupled with smoking.
All squamous cell carcinomas arising from the oropharynx, and all neck node metastases of unknown primary should undergo testing for HPV status. This is essential to adequately stage the tumour and adequately plan treatment. Due to the different biology of HPV positive and negative cancers, differentiating HPV status is also important for ongoing research to determine the best treatments.
Nasopharyngeal carcinomas, or neck node metastases possibly arising from the nasopharynx will also be tested for Ebstein Barr virus.
The tumor marker Cyfra 21-1 may be useful in diagnosing squamous cell carcinoma of the head and neck (SCCHN).
Adenocarcinoma
Adenocarcinoma is a cancer of the epithelial tissue that has glandular characteristics. Several head and neck cancers are adenocarcinomas (either of intestinal or non-intestinal cell types).
Prevention
Avoidance of risk factors (such as smoking and alcohol) is the single most effective form of prevention.
Regular dental examinations may identify pre-cancerous lesions in the oral cavity. While screening in the general population does not appear to be useful, screening high-risk groups by examination of the throat might be useful. Head and neck cancer is often curable if it is diagnosed early; however, outcomes are typically poor if it is diagnosed late.
When diagnosed early, oral, head, and neck cancers can be treated more easily, and the chances of survival increase tremendously. The HPV vaccine helps to prevent the development of HPV-related oropharyngeal cancer.
Management
Improvements in diagnosis and local management, as well as targeted therapy, have led to improvements in quality of life and survival for people with head and neck cancer.
After a histologic diagnosis has been established and tumor extent determined, such as with the use of PET-CT, the selection of appropriate treatment for a specific cancer depends on a complex array of variables, including tumor site, relative morbidity of various treatment options, concomitant health problems, social and logistic factors, previous primary tumors, and the person's preference. Treatment planning generally requires a multidisciplinary approach involving specialist surgeons, medical oncologists, and radiation oncologists.
Surgical resection and radiation therapy are the mainstays of treatment for most head and neck cancers and remain the standard of care in most cases. For small primary cancers without regional metastases (stage I or II), wide surgical excision alone or curative radiation therapy alone is used. For more extensive primary tumors or those with regional metastases (stage III or IV), planned combinations of pre- or postoperative radiation and complete surgical excision are generally used. More recently, as historical survival and control rates have been recognized as less than satisfactory, there has been an emphasis on the use of various induction or concomitant chemotherapy regimens.
Surgery
Surgery as a treatment is frequently used for most types of head and neck cancer. Usually, the goal is to remove the cancerous cells entirely. This can be particularly tricky if the cancer is near the larynx and can result in the person being unable to speak. Surgery is also commonly used to resect (remove) some or all of the cervical lymph nodes to prevent further spread of the disease. Transoral robotic surgery (TORS) is gaining popularity worldwide as the technology and training become more accessible. It now has an established role in the treatment of early stage oropharyngeal cancer. There is also a growing trend worldwide towards TORS for the surgical treatment of laryngeal and hypopharyngeal tumours.
CO2 laser surgery is also another form of treatment. Transoral laser microsurgery allows surgeons to remove tumors from the voice box with no external incisions. It also allows access to tumors that are not reachable with robotic surgery. During the surgery, the surgeon and pathologist work together to assess the adequacy of excision ("margin status"), minimizing the amount of normal tissue removed or damaged. This technique helps give the person as much speech and swallowing function as possible after surgery.
Radiation therapy
Radiation therapy is the most common form of treatment. There are different forms of radiation therapy, including 3D conformal radiation therapy, intensity-modulated radiation therapy, particle beam therapy, and brachytherapy, which are commonly used in the treatment of cancers of the head and neck. Most people with head and neck cancer who are treated in the United States and Europe are treated with intensity-modulated radiation therapy using high-energy photons. At higher doses, head and neck radiation is associated with thyroid dysfunction and pituitary axis dysfunction. Radiation therapy for head and neck cancers can also cause acute skin reactions of varying severity, which can be treated and managed with topically applied creams or specialist films.
Chemotherapy
Chemotherapy for throat cancer is not generally used to cure the cancer as such. Instead, it is used to provide an inhospitable environment for metastases so that they will not establish themselves in other parts of the body. Typical chemotherapy agents are a combination of paclitaxel and carboplatin. Cetuximab is also used in the treatment of throat cancer.
Docetaxel-based chemotherapy has shown a very good response in locally advanced head and neck cancer. Docetaxel is the only taxane approved by the FDA for head and neck cancer, in combination with cisplatin and fluorouracil for the induction treatment of inoperable, locally advanced head and neck cancer.
While not specifically a chemotherapy, amifostine is often administered intravenously by a chemotherapy clinic prior to IMRT radiotherapy sessions. Amifostine protects the gums and salivary glands from the effects of radiation.
There is no evidence that erythropoietin should be routinely given with radiotherapy.
Photodynamic therapy
Photodynamic therapy may have promise for treating mucosal dysplasia and small head and neck tumors. Amphinex is showing good results in early clinical trials for the treatment of advanced head and neck cancer.
Targeted therapy
Targeted therapy, according to the National Cancer Institute, is "a type of treatment that uses drugs or other substances, such as monoclonal antibodies, to identify and attack specific cancer cells without harming normal cells." Some targeted therapies used in head and neck cancers include cetuximab, bevacizumab, and erlotinib.
Cetuximab is used for treating people with advanced-stage cancer who cannot be treated with conventional chemotherapy (cisplatin). However, cetuximab's efficacy is still under investigation by researchers.
Gendicine is a gene therapy that employs an adenovirus to deliver the tumor suppressor gene p53 to cells. It was approved in China in 2003 for the treatment of head and neck cancer.
The mutational profiles of HPV+ and HPV- head and neck cancer have been reported, further demonstrating that they are fundamentally distinct diseases.
Immunotherapy
Immunotherapy is a type of treatment that activates the immune system to fight cancer. One type of immunotherapy, immune checkpoint blockade, binds to and blocks inhibitory signals on immune cells to release their anti-cancer activities.
In 2016, the FDA granted accelerated approval to pembrolizumab for the treatment of people with recurrent or metastatic head and neck cancer with disease progression on or after platinum-containing chemotherapy. Later that year, the FDA approved nivolumab for the treatment of recurrent or metastatic head and neck cancer with disease progression on or after platinum-based chemotherapy. In 2019, the FDA approved pembrolizumab for the first-line treatment of metastatic or unresectable recurrent head and neck cancer.
Treatment side effects
Depending on the treatment used, people with head and neck cancer may experience the following symptoms and treatment side effects:
Eating problems
Pain associated with lesions
Mucositis
Nephrotoxicity and ototoxicity
Xerostomia
Gastroesophageal reflux
Radiation-induced osteonecrosis of the jaw
Radiation-induced acute skin reactions
Psychosocial
Programs to support the emotional and social well-being of people who have been diagnosed with head and neck cancer may be offered. There is no clear evidence on the effectiveness of these interventions or any particular type of psychosocial program or length of time that is most helpful for those with head and neck cancer.
Prognosis
Although early-stage head and neck cancers (especially laryngeal and oral cavity) have high cure rates, up to 50% of people with head and neck cancer present with advanced disease.
Cure rates decrease in locally advanced cases, whose probability of cure is inversely related to tumor size and even more so to the extent of regional node involvement. HPV-associated oropharyngeal cancer has been shown to respond better to chemoradiation and, subsequently, have a better prognosis compared to non-associated HPV head and neck cancer.
Consensus panels in America (AJCC) and Europe (UICC) have established staging systems for head and neck cancers. These staging systems attempt to standardize clinical trial criteria for research studies and define prognostic categories of disease. Head and neck cancers are staged according to the TNM classification system, where T is the size and configuration of the tumor, N is the presence or absence of lymph node metastases, and M is the presence or absence of distant metastases. The T, N, and M characteristics are combined to produce a "stage" of the cancer, from I to IVB.
Problem of second primaries
Survival advantages provided by new treatment modalities have been undermined by the significant percentage of people cured of head and neck cancer who subsequently develop second primary tumors. The incidence of second primary tumors ranges in studies from 9%
to 23%
at 20 years. Second primary tumors are the major threat to long-term survival after successful therapy of early-stage head and neck cancer. Their high incidence results from the same carcinogenic exposure responsible for the initial primary process, called field cancerization.
Digestive system
Many people with head and neck cancer are also not able to eat sufficiently. A tumor may impair a person's ability to swallow and eat, and throat cancer may affect the digestive system. The difficulty in swallowing can cause a person to choke on their food in the early stages of digestion and interfere with the food's smooth travel down into the esophagus and beyond.
The treatments for throat cancer can also be harmful to the digestive system as well as other body systems. Radiation therapy can lead to nausea and vomiting, which can deprive the body of vital fluids (although these may be obtained through intravenous fluids if necessary). Frequent vomiting can lead to an electrolyte imbalance, which has serious consequences for the proper functioning of the heart. Frequent vomiting can also upset the balance of stomach acids, which has a negative impact on the digestive system, especially the lining of the stomach and esophagus.
Enteral feeding, a method that adds nutrients directly into a person's stomach using a nasogastric feeding tube or a gastrostomy tube, may be necessary for some people. Further research is required to determine the most effective method of enteral feeding to ensure that people undergoing radiotherapy or chemoradiation treatment are able to stay nourished during their treatment.
Mental health
Cancer in the head or neck may impact a person's mental well-being and can sometimes lead to social isolation. This largely results from a decreased ability or inability to eat, speak, or effectively communicate. Physical appearance is often altered by the cancer itself and/or as a consequence of treatment side effects. Psychological distress may occur, and feelings such as uncertainty and fear may arise. Some people may also have a changed physical appearance, differences in swallowing or breathing, and residual pain to manage.
Caregiver stress
Caregivers for people with head and neck cancer show higher rates of caregiver stress and poorer mental health compared to both the general population and those caring for non-head and neck cancer patients. The high symptom burden patients' experience necessitates complex caregiver roles, often requiring hospital staff training, which caregivers can find distressing when asked to do so for the first time. It is becoming increasingly apparent that caregivers (most often spouses, children, or close family members) might not be adequately informed about, prepared for, or trained for the tasks and roles they will encounter during the treatment and recovery phases of this unique patient population, which span both technical and emotional support. Of note, caregivers of patients who report lower quality of life demonstrate increased burden and fatigue that extend beyond the treatment phase.
Examples of technically difficult caregiver duties include tube feeding, oral suctioning, wound maintenance, medication delivery safe for tube feeding, and troubleshooting home medical equipment. If the cancer affects the mouth or larynx, caregivers must also find a way to effectively communicate among themselves and with their healthcare team. This is in addition to providing emotional support for the person undergoing cancer therapy.
Others
Like any cancer, metastasis affects many areas of the body as the cancer spreads from cell to cell and organ to organ. For example, if it spreads to the bone marrow, it will prevent the body from producing enough red blood cells and affect the proper functioning of the white blood cells and the body's immune system; spreading to the circulatory system will prevent oxygen from being transported to all the cells of the body; and throat cancer can throw the nervous system into chaos, making it unable to properly regulate and control the body.
Epidemiology
Globally, head and neck cancer accounts for 650,000 new cases of cancer and 330,000 deaths annually on average. In 2018, it was the seventh most common cancer worldwide, with 890,000 new cases documented and 450,000 people dying from the disease. The risk of developing head and neck cancer increases with age, especially after 50 years. Most people who do so are between 50 and 70 years old.
In North America and Europe, the tumors usually arise from the oral cavity, oropharynx, or larynx, whereas nasopharyngeal cancer is more common in the Mediterranean countries and in the Far East. In Southeast China and Taiwan, head and neck cancer, specifically nasopharyngeal cancer, is the most common cause of death in young men.
United States
In the United States, head and neck cancer makes up 3% of all cancer cases (averaging 53,000 new diagnoses per year) and 1.5% of cancer deaths. The 2017 worldwide figure cites head and neck cancers as representing 5.3% of all cancers (not including non-melanoma skin cancers). Notably, head and neck cancer secondary to chronic alcohol or tobacco use has been steadily declining as less of the population chronically smokes tobacco. However, HPV-associated oropharyngeal cancer is rising, particularly in younger people in westernized nations, which is thought to be reflective of changes in oral sexual practices, specifically with regard to the number of oral sexual partners. This increase since the 1970s has mostly affected wealthier nations and male populations. This is due to evidence suggesting that transmission rates of HPV from women to men are higher than from men to women, as women often have a higher immune response to infection.
In 2008, there were 22,900 cases of oral cavity cancer, 12,250 cases of laryngeal cancer, and 12,410 cases of pharyngeal cancer in the United States.
In 2002, 7,400 Americans were projected to die of these cancers.
More than 70% of throat cancers are at an advanced stage when discovered.
Men are 89% more likely than women to be diagnosed with these cancers and are almost twice as likely to die of them.
African Americans are disproportionately affected by head and neck cancer, with younger ages of incidence, increased mortality, and more advanced disease at presentation. Laryngeal cancer incidence is higher in African Americans relative to white, Asian, and Hispanic populations. There is a lower survival rate for similar tumor states in African Americans with head and neck cancer.
Research
Immunotherapy with immune checkpoint inhibitors is being investigated in head and neck cancers.
| Biology and health sciences | Cancer | Health |
729876 | https://en.wikipedia.org/wiki/Microreactor | Microreactor | A microreactor or microstructured reactor or microchannel reactor is a device in which chemical reactions take place in a confinement with typical lateral dimensions below 1 mm;
the most typical form of such confinement are microchannels. Microreactors are studied in the field of micro process engineering, together with other devices (such as micro heat exchangers) in which physical processes occur. The microreactor is usually a continuous flow reactor (contrast with/to a batch reactor). Microreactors can offer many advantages over conventional scale reactors, including improvements in energy efficiency, reaction speed and yield, safety, reliability, scalability, on-site/on-demand production, and a much finer degree of process control.
History
Gas-phase microreactors have a long history but those involving liquids started to appear in the late 1990s. One of the first microreactors with embedded high performance heat exchangers were made in the early 1990s by the Central Experimentation Department (Hauptabteilung Versuchstechnik, HVT) of Forschungszentrum Karlsruhe
in Germany, using mechanical micromachining techniques that were a spinoff from the manufacture of separation nozzles for uranium enrichment. As research on nuclear technology was drastically reduced in Germany, microstructured heat exchangers were investigated for their application in handling highly exothermic and dangerous chemical reactions. This new concept, known by names as microreaction technology or micro process engineering, was further developed by various research institutions. An early example from 1997 involved that of azo couplings in a pyrex reactor with channel dimensions 90 micrometres deep and 190 micrometres wide.
Benefits
Using microreactors is somewhat different from using a glass vessel. These reactors may be a valuable tool in the hands of an experienced chemist or reaction engineer:
Microreactors typically have heat exchange coefficients of at least 1 megawatt per cubic meter per kelvin, up to 500 MW m−3 K−1 vs. a few kilowatts in conventional glassware (1 L flask ~10 kW m−3 K−1). Thus, microreactors can remove heat much more efficiently than vessels and even critical reactions such as nitrations can be performed safely at high temperatures. Hot spot temperatures as well as the duration of high temperature exposition due to exothermicity decreases remarkably. Thus, microreactors may allow better kinetic investigations, because local temperature gradients affecting reaction rates are much smaller than in any batch vessel. Heating and cooling a microreactor is also much quicker and operating temperatures can be as low as −100 °C. As a result of the superior heat transfer, reaction temperatures may be much higher than in conventional batch-reactors. Many low temperature reactions as organo-metal chemistry can be performed in microreactors at temperatures of −10 °C rather than −50 °C to −78 °C as in laboratory glassware equipment.
Microreactors are normally operated continuously. This allows the subsequent processing of unstable intermediates and avoids typical batch workup delays. Especially low temperature chemistry with reaction times in the millisecond to second range are no longer stored for hours until dosing of reagents is finished and the next reaction step may be performed. This rapid work up avoids decay of precious intermediates and often allows better selectivities.
Continuous operation and mixing causes a very different concentration profile when compared with a batch process. In a batch, reagent A is filled in and reagent B is slowly added. Thus, B encounters initially a high excess of A. In a microreactor, A and B are mixed nearly instantly and B won't be exposed to a large excess of A. This may be an advantage or disadvantage depending on the reaction mechanism - it is important to be aware of such different concentration profiles.
Although a bench-top microreactor can synthesize chemicals only in small quantities, scale-up to industrial volumes is simply a process of multiplying the number of microchannels. In contrast, batch processes too often perform well on R&D bench-top level but fail at batch pilot plant level.
Pressurisation of materials within microreactors (and associated components) is generally easier than with traditional batch reactors. This allows reactions to be increased in rate by raising the temperature beyond the boiling point of the solvent. This, although typical Arrhenius behaviour, is more easily facilitated in microreactors and should be considered a key advantage. Pressurisation may also allow dissolution of reactant gasses within the flow stream.
Challenges
Although there have been reactors made for handling particles, microreactors generally do not tolerate particulates well, often clogging. Clogging has been identified by a number of researchers as the biggest hurdle for microreactors being widely accepted as a beneficial alternative to batch reactors. So far, the so-called microjetreactor is free of clogging by precipitating products. Gas evolved may also shorten the residence time of reagents as volume is not constant during the reaction. This may be prevented by application of pressure.
Mechanical pumping may generate a pulsating flow which can be disadvantageous. Much work has been devoted to development of pumps with low pulsation. A continuous flow solution is electroosmotic flow (EOF).
The logistics issue and enhanced pressure drop across microreactor limits its applicability in large-scale production units. However, clean solutions are handled well in microreactors.
The scale-up of production rates and leakage are quite challenging in case of microreactor. Recently, so called nanoparticle immobilized reactors are developed to solve logistics and scale-up issues, associated with the microreactors.
Typically, reactions performing very well in a microreactor encounter many problems in vessels, especially when scaling up. Often, the high area to volume ratio and the uniform residence time cannot easily be scaled.
Corrosion imposes a bigger issue in microreactors because area to volume ratio is high. Degradation of few μm may go unnoticed in conventional vessels. As typical inner dimensions of channels are in the same order of magnitude, characteristics may be altered significantly.
T reactors
One of the simplest forms of a microreactor is a 'T' reactor. A 'T' shape is etched into a plate with a depth that may be 40 micrometres and a width of 100 micrometres: the etched path is turned into a tube by sealing a flat plate over the top of the etched groove. The cover plate has three holes that align to the top-left, top-right, and bottom of the 'T' so that fluids can be added and removed. A solution of reagent 'A' is pumped into the top left of the 'T' and solution 'B' is pumped into the top right of the 'T'. If the pumping rate is the same, the components meet at the top of the vertical part of the 'T' and begin to mix and react as they go down the trunk of the 'T'. A solution of product is removed at the base of the 'T'.
Applications
Synthesis
Microreactors can be used to synthesise material more effectively than current batch techniques allow. The benefits here are primarily enabled by the mass transfer, thermodynamics, and high surface area to volume ratio environment as well as engineering advantages in handling unstable intermediates. Microreactors are applied in combination with photochemistry, electrosynthesis, multicomponent reactions and polymerization (for example that of butyl acrylate). It can involve liquid-liquid systems but also solid-liquid systems with for example the channel walls coated with a heterogeneous catalyst. Synthesis is also combined with online purification of the product. Following green chemistry principles, microreactors can be used to synthesize and purify extremely reactive Organometallic Compounds for ALD and CVD applications, with improved safety in operations and higher purity products.
In microreactor studies a Knoevenagel condensation was performed with the channel coated with a zeolite catalyst layer which also serves to remove water generated in the reaction. The same reaction was performed in a microreactor covered by polymer brushes.
A Suzuki reaction was examined in another study with a palladium catalyst confined in a polymer network of polyacrylamide and a triarylphosphine formed by interfacial polymerization:
The combustion of propane was demonstrated to occur at temperatures as low as 300 °C in a microchannel setup filled up with an aluminum oxide lattice coated with a platinum / molybdenum catalyst:
Enzyme catalyzed polymer synthesis
Enzymes immobilized on solid supports are increasingly used for greener, more sustainable chemical transformation processes. > enabled to perform heterogeneous reactions in continuous mode, in organic media, and at elevated temperatures. Using microreactors, enabled faster polymerization and higher molecular mass compared to using batch reactors. It is evident that similar microreactor based platforms can readily be extended to other enzyme-based systems, for example, high-throughput screening of new enzymes and to precision measurements of new processes where continuous flow mode is preferred. This is the first reported demonstration of a solid supported enzyme-catalyzed polymerization reaction in continuous mode.
Analysis
Microreactors can also enable experiments to be performed at a far lower scale and far higher experimental rates than currently possible in batch production, while not collecting the physical experimental output. The benefits here are primarily derived from the low operating scale, and the integration of the required sensor technologies to allow high quality understanding of an experiment. The integration of the required synthesis, purification and analytical capabilities is impractical when operating outside of a microfluidic context.
NMR
Researchers at the Radboud University Nijmegen and Twente University, the Netherlands, have developed a microfluidic high-resolution NMR flow probe. They have shown a model reaction being followed in real-time. The combination of the uncompromised (sub-Hz) resolution and a low sample volume can prove to be a valuable tool for flow chemistry.
Infrared spectroscopy
Mettler Toledo and Bruker Optics offer dedicated equipment for monitoring, with attenuated total reflectance spectrometry (ATR spectrometry) in microreaction setups. The former has been demonstrated for reaction monitoring. The latter has been successfully used for reaction monitoring and determining dispersion characteristics of a microreactor.
Academic research
Microreactors, and more generally, micro process engineering, are the subject of worldwide academic research. A prominent recurring conference is IMRET, the International Conference on Microreaction Technology. Microreactors and micro process engineering have also been featured in dedicated sessions of other conferences, such as the Annual Meeting of the American Institute of Chemical Engineers (AIChE), or the International Symposia on Chemical Reaction Engineering (ISCRE). Research is now also conducted at various academic institutions around the world, e.g. at the Massachusetts Institute of Technology (MIT) in Cambridge, Massachusetts, University of Illinois Urbana-Champaign, Oregon State University in Corvallis, Oregon, at University of California, Berkeley in Berkeley, California in the United States, at the EPFL in Lausanne, Switzerland, at Eindhoven University of Technology in Eindhoven, at Radboud University Nijmegen in Nijmegen, Netherlands and at the LIPHT of Université de Strasbourg in Strasbourg and LGPC of the University of Lyon, CPE Lyon, France and at KU Leuven, Belgium.
Market structure
The market for microreactors can be segmented based on customer objectives, into turnkey, modular, and bespoke systems.
Turnkey (ready to run) systems are being used where the application environment stands to benefit from new chemical synthesis schemes, enhanced investigational throughput of up to approximately 10 - 100 experiments per day (depends on reaction time) and reaction subsystem, and actual synthesis conduct at scales ranging from 10 milligrams per experiment to triple digit tons per year (continuous operation of a reactor battery).
Modular (open) systems are serving the niche for investigations on continuous process engineering lay-outs, where a measurable process advantage over the use of standardized equipment is anticipated by chemical engineers. Multiple process lay-outs can be rapidly assembled and chemical process results obtained on a scale ranging from several grams per experiment up to approximately 100 kg at a moderate number of experiments per day (3-15). A secondary transfer of engineering findings in the context of a plant engineering exercise (scale-out) then provides target capacity of typically single product dedicated plants. This mimics the success of engineering contractors for the petro-chemical process industry.
With dedicated developments, manufacturers of microstructured components are mostly commercial development partners to scientists in search of novel synthesis technologies. Such development partners typically excel in the set-up of comprehensive investigation and supply schemes, to model a desired contacting pattern or spatial arrangement of matter. To do so they predominantly offer information from proprietary integrated modeling systems that combine computational fluid dynamics with thermokinetic modelling. Moreover, as a rule, such development partners establish the overall application analytics to the point where the critical initial hypothesis can be validated and further confined.
| Physical sciences | Chemical engineering | Chemistry |
729998 | https://en.wikipedia.org/wiki/Zinc%20sulfide | Zinc sulfide | Zinc sulfide (or zinc sulphide) is an inorganic compound with the chemical formula of ZnS. This is the main form of zinc found in nature, where it mainly occurs as the mineral sphalerite. Although this mineral is usually black because of various impurities, the pure material is white, and it is widely used as a pigment. In its dense synthetic form, zinc sulfide can be transparent, and it is used as a window for visible optics and infrared optics.
Structure
ZnS exists in two main crystalline forms. This dualism is an example of polymorphism. In each form, the coordination geometry at Zn and S is tetrahedral. The more stable cubic form is known also as zinc blende or sphalerite. The hexagonal form is known as the mineral wurtzite, although it also can be produced synthetically. The transition from the sphalerite form to the wurtzite form occurs at around 1020 °C.
Applications
Luminescent material
Zinc sulfide, with addition of a few ppm of a suitable activator, exhibits strong phosphorescence. The phenomenon was described by Nikola Tesla in 1893, and is currently used in many applications, from cathode-ray tubes through X-ray screens to glow in the dark products. When silver is used as activator, the resulting color is bright blue, with maximum at 450 nanometers. Using manganese yields an orange-red color at around 590 nanometers. Copper gives a longer glow, and it has the familiar greenish glow-in-the-dark. Copper-doped zinc sulfide ("ZnS plus Cu") is used also in electroluminescent panels. It also exhibits phosphorescence due to impurities on illumination with blue or ultraviolet light.
Optical material
Zinc sulfide is also used as an infrared optical material, transmitting from visible wavelengths to just over 12 micrometers. It can be used planar as an optical window or shaped into a lens. It is made as microcrystalline sheets by the synthesis from hydrogen sulfide gas and zinc vapour, and this is sold as FLIR-grade (Forward Looking Infrared), where the zinc sulfide is in a milky-yellow, opaque form. This material when hot isostatically pressed (HIPed) can be converted to a water-clear form known as Cleartran (trademark). Early commercial forms were marketed as Irtran-2 but this designation is now obsolete.
Pigment
Zinc sulfide is a common pigment, sometimes called sachtolith. When combined with barium sulfate, zinc sulfide forms lithopone.
Catalyst
Fine ZnS powder is an efficient photocatalyst, which produces hydrogen gas from water upon illumination. Sulfur vacancies can be introduced in ZnS during its synthesis; this gradually turns the white-yellowish ZnS into a brown powder, and boosts the photocatalytic activity through enhanced light absorption.
Semiconductor properties
Both sphalerite and wurtzite are intrinsic, wide-bandgap semiconductors. These are prototypical II-VI semiconductors, and they adopt structures related to many of the other semiconductors, such as gallium arsenide. The cubic form of ZnS has a band gap of about 3.54 electron volts at 300 kelvins, but the hexagonal form has a band gap of about 3.91 electron volts. ZnS can be doped as either an n-type semiconductor or a p-type semiconductor.
History
The phosphorescence of ZnS was first reported by the French chemist Théodore Sidot in 1866. His findings were presented by A. E. Becquerel, who was renowned for the research on luminescence. ZnS was used by Ernest Rutherford and others in the early years of nuclear physics as a scintillation detector, because it emits light upon excitation by x-rays or electron beam, making it useful for X-ray screens and cathode-ray tubes. This property made zinc sulfide useful in the dials of radium watches.
Production
Zinc sulfide is usually produced from waste materials from other applications. Typical sources include smelter, slag, and pickle liquors. As an example, the synthesis of ammonia from methane requires a priori removal of hydrogen sulfide impurities in the natural gas, for which zinc oxide is used. This scavenging produces zinc sulfide:
ZnO + H2S → ZnS + H2O
Laboratory preparation
Crude zinc sulfide can be produced by igniting a mixture of zinc and sulfur. More conventionally, ZnS is prepared by treating a mildly acidic solution of Zn2+ salts with H2S:
Zn2+ + S2− → ZnS
This reaction is the basis of a gravimetric analysis for zinc.
| Physical sciences | Sulfide salts | Chemistry |
730219 | https://en.wikipedia.org/wiki/Bathing | Bathing | Bathing is the immersion of the body, wholly or partially, usually in water, but often in another medium such as hot air. It is most commonly practised as part of personal cleansing, and less frequently for relaxation or as a leisure activity. Cleansing the body may be solely a component of personal hygiene, but is also a spiritual part of some religious rituals. Bathing is also sometimes used medically or therapeutically, as in hydrotherapy, ice baths, or the mud bath.
People bathe in water at temperatures ranging from very cold to very hot, or in appropriately heated air, according to custom or purpose.
Where indoor heated water is available, people bathe more or less daily, at comfortable temperatures, in a private bathtub or shower. Communal bathing, such as that in hammams, sauna, banya, Victorian Turkish baths, and sentō, fulfils the same purpose, in addition to its often having a social function.
Ritual religious bathing is sometimes referred to as immersion. This can be required after sexual intercourse or menstruation (Islam and Judaism), or as baptism (Christianity).
By analogy, the term "bathing" is also applied to relaxing activities in which the participant "bathes" in the rays of the sun (sunbathing) or in outdoor bodies of water, such as in sea bathing or wild swimming.
Although there is sometimes overlap, as in sea bathing, most bathing is usually treated as distinct from more active recreations like swimming.
History
Ancient world
Bathing in Ancient China may be traced back to the Shang Dynasty, 3000 years ago (1600–1046 BCE). Archaeological findings from the Yinxu ruins show a cauldron to boil water, smaller cauldrons to draw out the water to be poured into a basin, skin scrapers to remove dirt and dead skin. 2300 year old lavish imperial bathrooms with exquisite tiles and a sewage system can be seen in Xi'an. Bathing grew in importance in the Han Dynasty (202 BC–AD 220) where officials were allowed to take a day's leave for bathing at home every five days, and bathing became the reason for a bank holiday for the first time.
An accountable daily ritual of bathing can be traced to the ancient Indians. They used elaborate practices for personal hygiene with three daily baths and washing. These are recorded in the works called grihya sutras which date back to 500 BCE and are in practice today in some communities. In Hinduism, “Prataha Kaal” (the onset of day) or “Brahma Muhoortham” begins with the 4 am “snanam” or bath, and was considered extremely auspicious in ancient times.
Ancient Greece utilized small bathtubs, wash basins, and foot baths for personal cleanliness. The earliest findings of baths date from the mid-2nd millennium BC in the palace complex at Knossos, Crete, and the luxurious alabaster bathtubs excavated in Akrotiri, Santorini. A word for bathtub, (), occurs eleven times in Homer. As a legitimate Mycenaean word (a-sa-mi-to) for a kind of vessel that could be found in any Mycenaean palace, this Linear B term derives from an Aegean suffix -inth- being appended to an Akkadian loan word with the root namsû ('washbowl', 'washing tub'). This luxurious item of the Mycenaean palace culture, therefore, was clearly borrowed from the Near East. Later Greeks established public baths and showers within gymnasiums for relaxation and personal hygiene. The word gymnasium (γυμνάσιον) comes from the Greek word gymnos (γυμνός), meaning "naked".
Ancient Rome developed a network of aqueducts to supply water to all large towns and population centers and had indoor plumbing, with pipes that terminated in homes and at public wells and fountains. The Roman public baths were called thermae. The thermae were not simply baths, but important public works that provided facilities for many kinds of physical exercise and ablutions, with cold, warm, and hot baths, rooms for instruction and debate, and usually one Greek and one Latin library. They also represented an important moment of socialization and exchange between the members of the community. They were provided for the public by a benefactor, usually the Emperor. Other empires of the time did not show such an affinity for public works, but this Roman practice spread their culture to places where there may have been more resistance to foreign mores. Unusually for the time, the thermae were not class-stratified, being available to all for no charge or a small fee. With the fall of the Roman Empire, the aqueduct system fell into disrepair and disuse. But even before that, during the Christianization of the Empire, changing ideas about public morals led the baths into disfavor.
Medieval Japan
Before the 7th century, the Japanese were likely to have bathed in the many springs in the open, as there is no evidence of closed rooms. In the 6th to 8th centuries (in the Asuka and Nara periods) the Japanese absorbed the religion of Buddhism from China, which had a strong impact on the culture of the entire country. Buddhist temples traditionally included a bathhouse (yuya) for the monks. Due to the principle of purity espoused by Buddhism these baths were eventually opened to the public. Only the wealthy had private baths.
The first public bathhouse was mentioned in 1266. In Edo (modern Tokyo), the first sentō was established in 1591. The early steam baths were called iwaburo ( "rock pools") or kamaburo ( "furnace baths"). These were built into natural caves or stone vaults. In iwaburo along the coast, the rocks were heated by burning wood, then sea water was poured over the rocks, producing steam. The entrances to these "bath houses" were very small, possibly to slow the escape of the heat and steam. There were no windows, so it was very dark inside and the user constantly coughed or cleared their throats in order to signal to new entrants which seats were already occupied. The darkness could be also used to cover sexual contact. Because there was no gender distinction, these baths came into disrepute. They were finally abolished in 1870 on hygienic and moral grounds. Author John Gallagher says bathing "was segregated in the 1870s as a concession to outraged Western tourists".
At the beginning of the Edo period (1603–1868) there were two different types of baths. In Edo, hot-water baths (' ) were common, while in Osaka, steam baths ( ) were common. At that time shared bathrooms for men and women were the rule. These bathhouses were very popular, especially for men. "Bathing girls" ( ) were employed to scrub the guests' backs and wash their hair, etc. In 1841, the employment of yuna was generally prohibited, as well as mixed bathing. The segregation of the sexes, however, was often ignored by operators of bathhouses, or areas for men and women were separated only by a symbolic line. Today, sento baths have separate rooms for men and women.
Mesoamerica
Spanish chronicles describe the bathing habits of the peoples of Mesoamerica during and after the conquest.
Bernal Díaz del Castillo describes Moctezuma (the Mexica, or Aztec, king at the arrival of Cortés) in his Historia verdadera de la conquista de la Nueva España as being "...Very neat and cleanly, bathing every day each afternoon...".
Bathing was not restricted to the elite, but was practised by all people; the chronicler Tomás López Medel wrote after a journey to Central America that "Bathing and the custom of washing oneself is so quotidian (common) amongst the Indians, both of cold and hot lands, as is eating, and this is done in fountains and rivers and other water to which they have access, without anything other than pure water..."
The Mesoamerican bath, known as temazcal in Spanish, from the Nahuatl word temazcalli, a compound of temaz ("steam") and calli ("house"), consists of a room, often in the form of a small dome, with an exterior firebox known as texictle (teʃict͜ɬe) that heats a small portion of the room's wall made of volcanic rocks; after this wall has been heated, water is poured on it to produce steam, an action known as tlasas. As the steam accumulates in the upper part of the room a person in charge uses a bough to direct the steam to the bathers who are lying on the ground, with which he later gives them a massage, then the bathers scrub themselves with a small flat river stone and finally the person in charge introduces buckets with water with soap and grass used to rinse. This bath had also ritual importance, and was vinculated to the goddess Toci; it is also therapeutic when medicinal herbs are used in the water for the tlasas. It is still used in Mexico.
Medieval and early-modern Europe
Christianity has always placed a strong emphasis on hygiene. Despite the denunciation of the mixed bathing style of Roman pools by early Christian clergy, as well as the pagan custom of women bathing naked in front of men, this did not stop the Church from urging its followers to go to public baths for bathing, which contributed to hygiene and good health according to the Church Fathers, Clement of Alexandria and Tertullian. The Church also built public bathing facilities that were separate for both sexes near monasteries and pilgrimage sites; also, the popes situated baths within church basilicas and monasteries since the early Middle Ages. Pope Gregory the Great urged his followers on value of bathing as a bodily need.
Great bathhouses were built in Byzantine centers such as Constantinople and Antioch, and the popes allocated to the Romans bathing through diaconia, or private Lateran baths, or even a myriad of monastic bath houses functioning in eighth and ninth centuries. The Popes maintained their baths in their residences which described by scholar Paolo Squatriti as "luxurious baths", and bath houses including hot baths incorporated into Christian Church buildings or those of monasteries, which known as "charity baths" because they served both the clerics and needy poor people. Public bathing was common in larger towns and cities such as Paris, Regensburg and Naples. The Catholic religious orders of the Augustinians and Benedictines had rules for ritual purification, and inspired by Benedict of Nursia encouragement for the practice of therapeutic bathing; Benedictine monks played a role in the development and promotion of spas. Protestantism also played a prominent role in the development of the British spas.
In the Middle Ages, bathing commonly took place in public bathhouses. Public baths were also havens for prostitution, which created some opposition to them. Rich people bathed at home, most likely in their bedroom, as "bath" rooms were not common. Bathing was done in large, wooden tubs with a linen cloth laid in it to protect the bather from splinters. Additionally, during the Renaissance and Protestant Reformation, the quality and condition of the clothing (as opposed to the actual cleanliness of the body itself) were thought to reflect the soul of an individual. Clean clothing also reflected one's social status; clothes made the man or woman.
Due to Black Death plague, introduced from Asia to Europe, public baths were closed to avoid contagion. In the sixteenth century, the popularity of public bathhouses in Europe sharply declined, perhaps due to the new plague of syphilis which made sexual promiscuity more risky, or stronger religious prohibitions on nudity surrounding the Protestant Reformation. Some Europeans came to believe the false idea that bathing or steaming would open pores to disease.
Modern era
Therapeutic bathing
Public opinion about bathing began to shift in the middle and late 18th century, when writers argued that frequent bathing might lead to better health. Two English works on the medical uses of water were published in the 18th century that inaugurated the new fashion for therapeutic bathing. One of these was by Sir John Floyer, a physician of Lichfield, who, struck by the remedial use of certain springs by the neighbouring peasantry, investigated the history of cold bathing and published a book on the subject in 1702.
The book ran through six editions within a few years and the translation of this book into German was largely drawn upon by Dr J. S. Hahn of Silesia as the basis for his book called On the Healing Virtues of Cold Water, Inwardly and Outwardly Applied, as Proved by Experience, published in 1738.
The other work was a 1797 publication by Dr James Currie of Liverpool on the use of hot and cold water in the treatment of fever and other illness, with a fourth edition published not long before his death in 1805. It was also translated into German by Michaelis (1801) and Hegewisch (1807). It was highly popular and first placed the subject on a scientific basis. Hahn's writings had meanwhile created much enthusiasm among his countrymen, societies having been everywhere formed to promote the medicinal and dietetic use of water; in 1804 Professor E.F.C. Oertel of Anspach republished them and quickened the popular movement by the unqualified commendation of water drinking as a remedy for all diseases.
A popular revival followed the application of hydrotherapy around 1829, by Vincenz Priessnitz, a peasant farmer in Gräfenberg, then part of the Austrian Empire.
This revival was continued by a Bavarian priest, Sebastian Kneipp (1821–1897), "an able and enthusiastic follower" of Priessnitz, "whose work he took up where Priessnitz left it", after he read a treatise on the cold water cure. In Wörishofen (south Germany), Kneipp developed the systematic and controlled application of hydrotherapy for the support of medical treatment that was delivered only by doctors at that time. Kneipp's own book My Water Cure was published in 1886 with many subsequent editions, and translated into many languages.
Captain R. T. Claridge was responsible for introducing and promoting hydropathy in Britain, first in London in 1842, then with lecture tours in Ireland and Scotland in 1843. His 10-week tour in Ireland included Limerick, Cork, Wexford, Dublin and Belfast, over June, July and August 1843, with two subsequent lectures in Glasgow.
The acceptance of germ theory in the late 1800s provided scientific reasons for frequent bathing.
Public baths and wash-houses
Large public baths such as those found in the ancient world and the Ottoman Empire were revived during the 19th century. The first modern public baths were opened in Liverpool in 1829. The first known warm fresh-water public wash house was opened in May 1842.
The popularity of wash-houses was spurred by the newspaper interest in Kitty Wilkinson, an Irish immigrant "wife of a labourer" who became known as the Saint of the Slums. In 1832, during a cholera epidemic, Wilkinson took the initiative to offer the use of her house and yard to neighbours to wash their clothes, at a charge of a penny per week, and showed them how to use a chloride of lime (bleach) to get them clean. She was supported by the District Provident Society and William Rathbone. In 1842, Wilkinson was appointed baths superintendent.
In Birmingham, around ten private baths were available in the 1830s. Whilst the dimensions of the baths were small, they provided a range of services. A major proprietor of bath houses in Birmingham was a Mr. Monro who had had premises in Lady Well and Snow Hill. Private baths were advertised as having healing qualities and being able to cure people of diabetes, gout and all skin diseases, amongst others. On 19 November 1844, it was decided that the working class members of society should have the opportunity to access baths, in an attempt to address the health problems of the public. On 22 April and 23 April 1845, two lectures were delivered in the town hall urging the provision of public baths in Birmingham and other towns and cities.
After a period of campaigning by many committees, the Public Baths and Wash-houses Act received royal assent on 26 August 1846. The act empowered local authorities across the country to incur expenditure in constructing public swimming baths out of its own funds.
The first London public baths was opened at Goulston Square, Whitechapel, in 1847 with the Prince consort laying the foundation stone.
Soap promoted for personal cleanliness
By the mid-19th century, the English urbanised middle classes had formed an ideology of cleanliness that ranked alongside typical Victorian concepts, such as Christianity, respectability and social progress. The cleanliness of the individual became associated with his or her moral and social standing within the community and domestic life became increasingly regulated by concerns regarding the presentation of domestic sobriety and cleanliness.
The industry of soapmaking began on a small scale in the 1780s, with the establishment of a soap manufactory at Tipton by James Keir and the marketing of high-quality, transparent soap in 1789 by Andrew Pears of London. In 1807, Pears found a way of removing the impurities and refining the base soap before adding the delicate perfume of garden flowers, founding Pears soap. It was in the mid-19th century, though, that the large-scale consumption of soap by the middle classes, anxious to prove their social standing, drove forward the mass production and marketing of soap.
William Gossage produced low-priced, good-quality soap from the 1850s. William Hesketh Lever and his brother, James, bought a small soap works in Warrington in 1886 and founded what is still one of the largest soap businesses, formerly called Lever Brothers and now called Unilever. These soap businesses were among the first to employ large-scale advertising campaigns. In 1882, English actress and socialite Lillie Langtry became the poster-girl for Pears soap, and thus the first celebrity to endorse a commercial product.
Before the late 19th century, water to individual places of residence was rare. Many countries in Europe developed a water collection and distribution network. London water supply infrastructure developed through major 19th-century treatment works built in response to cholera threats, to modern large-scale reservoirs. By the end of the century, private baths with running hot water were increasingly common in affluent homes in America and Britain.
At the beginning of the 20th century, a weekly Saturday night bath had become common custom for most of the population. A half day's work on Saturday for factory workers allowed them some leisure to prepare for the Sunday day of rest. The half day off allowed time for the considerable labor of drawing, carrying, and heating water, filling the bath and then afterward emptying it. To economize, bath water was shared by all family members. Indoor plumbing became more common in the 20th century and commercial advertising campaigns pushing new bath products began to influence public ideas about cleanliness, promoting the idea of a daily shower or bath.
In the twenty-first century challenges to the need for soap to effect such everyday cleanliness and whether soap is needed to avoid body odor, appeared in media.
Hot-air baths
Hammam
A hammam is a type of steam bath or a place of public bathing associated with the Islamic world. It is a prominent feature in the culture of the Muslim world and was inherited from the model of the Roman thermae. Muslim bathhouses or hammams were historically found across the Middle East, North Africa, al-Andalus (Islamic Spain and Portugal), Central Asia, the Indian subcontinent, and in Southeastern Europe under Ottoman rule.
In Islamic cultures the significance of the hammam was both religious and civic: it provided for the needs of ritual ablutions but also provided for general hygiene in an era before private plumbing and served other social functions such as offering a gendered meeting place for men and for women. Archeological remains attest to the existence of bathhouses in the Islamic world as early as the Umayyad period (7th–8th centuries) and their importance has persisted up to modern times. Their architecture evolved from the layout of Roman and Greek bathhouses and featured a regular sequence of rooms: an undressing room, a cold room, a warm room, and a hot room. Heat was produced by furnaces which provided hot water and steam, while smoke and hot air was channeled through conduits under the floor.
In a modern hammam visitors undress themselves, while retaining some sort of modesty garment or loincloth, and proceed into progressively hotter rooms, inducing perspiration. They are then usually washed by male or female staff (matching the gender of the visitor) with the use of soap and vigorous rubbing, before ending by washing themselves in warm water. Unlike in Roman or Greek baths, bathers usually wash themselves with running water instead of immersing themselves in standing water since this is a requirement of Islam, though immersion in a pool used to be customary in the hammams of some regions such as Iran. While hammams everywhere generally operate in fairly similar ways, there are some regional differences both in usage and architecture.
Victorian Turkish baths
Victorian Turkish baths (inspired by the traditional Islamic bathhouse—the hammam—itself an adaptation of the ancient Roman baths) were introduced to Britain by David Urquhart, diplomat and sometime Member of Parliament for Stafford. He wanted, for political and personal reasons, to popularize Turkish culture in Britain. In 1850 he wrote The Pillars of Hercules, a book about his travels in 1848 through Spain and Morocco. He described the vaporous hot-air baths (little-changed since Roman times) which he visited, both there and in the Ottoman Empire. In 1856 Dr Richard Barter read Urquhart's book and worked with him to construct such a bath, intending to use it at his hydropathic establishment at St Ann(e)'s Hill, near Blarney, County Cork, Ireland. Barter realised that the human body could tolerate the more therapeutically effective higher temperatures in hot air which was dry rather than steamy. After a number of unsuccessful attempts, he opened the first modern bath of this type in 1856. He called it the "Improved" Turkish or Irish bath, now better known as the Victorian Turkish bath.
The following year, the first public bath of its type to be built in mainland Britain since Roman times was opened in Manchester, and the idea spread rapidly. It reached London in July 1860, when Roger Evans, a member of one of Urquhart's Foreign Affairs Committees, opened a Turkish bath at 5 Bell Street, near Marble Arch. During the following 150 years, over 700 Turkish baths opened in the British Isles, including those built by municipal authorities as part of swimming pool complexes. It was claimed by Durham Dunlop (and many others) that hot-air bathing was a more effective body-cleanser than water, while Richard Metcalfe meticulously calculated that it would be more cost-effective for local authorities to provide hot-air baths in place of slipper baths.
Turkish baths opened in other parts of the British Empire. Dr. John Le Gay Brereton opened one in Sydney, Australia in 1859, Canada had one by 1869, and the first in New Zealand was opened in 1874. Urquhart's influence was also felt outside the Empire when in 1861, Dr Charles H. Shepard opened the first Turkish baths in the United States at 63 Columbia Street, Brooklyn Heights, New York, most probably on 3 October 1863.
Purpose
One purpose of bathing is for personal hygiene. It is a means of achieving cleanliness by washing away dead skin cells, dirt, and soil and as a preventative measure to reduce the incidence and spread of disease. It also may reduce body odors, however, some people note that may not be so necessary as commonly thought.
Bathing creates a feeling of well-being and the physical appearance of cleanliness.
Bathing may also be practised for religious ritual or therapeutic purposes or as a recreational activity. Bathing may be used to cool or to warm the body of an individual.
Therapeutic use of bathing includes hydrotherapy, healing, rehabilitation from injury or addiction, and relaxation.
The use of a bath in religious ritual or ceremonial rites include immersion during baptism in Christianity and to achieve a state of ritual cleanliness in a mikvah in Judaism. It is referred to as Ghusl in Arabic to attain ceremonial purity (Taahir) in Islam. All major religions place an emphasis on ceremonial purity, and bathing is one of the primary means of attaining outward purity. In Hindu households, any acts of defilement are countered by undergoing a bath and Hindus also immerse in Sarovar as part of religious rites. In the Sikh religion, there is a place at Golden Temple where the leprosy of Rajni's husband was cured by immersion into the holy sacred pool, and many pilgrims bathe in the sacred pool believing it will cure their illness as well.
Types of baths
Where bathing is for personal hygiene, bathing in a bathtub or shower is the most common form of bathing in Western, and many Eastern, countries. People most commonly bathe in their home or use a private bath in a public bathhouse. In some societies, bathing can take place in rivers, creeks, lakes or water holes, or any other place where there is an adequate pool of water. The quality of water used for bathing purposes varies considerably. Normally bathing involves use of soap or a soap-like substance, such as shower gel. In southern India people more commonly use aromatic oil and other home-made body scrubs.
Bathing in public can also provide occasions for social interaction, such as in Turkish, banya, sauna, hammams, or whirlpool baths.
Sponge bath
When water is in short supply or a person is not fit to have a standing bath, a wet cloth or sponge can be used, or the person can wash by splashing water over their body. A sponge bath is usually conducted in hospitals, which involves one person washing another with a sponge, while the person being washed remains lying in bed.
Ladling water from a container
This method involves using a small container to scoop water out of a large container and pour water over the body, in such a way that this water does not go back into the large container.
In Indonesia and Malaysia, this is a traditional method referred to as mandi.
In the Indonesian language, mandi is the verb for this process; bak mandi is the large container, and kamar mandi is the place in which this is done. Travel guides often use the word mandi on its own or in various ways such as for the large container and for the process of bathing.
In the Philippines, timba (pail) and tabo (dipper) are two essentials in every bathroom.
Bathing babies
Babies can be washed in a kitchen sink or a small plastic baby bath, instead of using a standard bath which offers little control of the infant's movements and requires the parent to lean awkwardly or kneel. Bathing infants too often has been linked to the development of asthma or severe eczema according to some researchers, including Michael Welch, chair of the American Academy of Pediatrics' section on allergy and immunology. A safe temperature for the bathwater is generally held to be .
Japanese bathing culture
Private baths
Although baths at home had been prevalent to some extent since the Edo period (1603-1867), the common people usually went to public bathhouses during the first half of the Showa period (1926-1989), and only wealthy families had their own bathrooms. Home baths became commonplace since the period of rapid economic growth after World War II. Bath water in Japan is much hotter than what is usual in Central Europe. The temperature is usually well above . In medical literature, is considered bearable. The heat is considered a prerequisite for complete relaxation. The custom is to thoroughly clean oneself with soap and rinse before entering the tub, so as not to contaminate the bath water.
Public baths
In public baths, there is a distinction between public baths with natural hot springs (called , meaning 'hot'), and those without natural hot springs (known as ). Since Japan is located in a volcanically active region, there are many hot springs, of which about 2,000 are swimming pools. Most are in the open countryside, but they are also found in cities. In Tokyo, for example, there are about 25 baths. Onsen are similar to Western-style spas in their therapeutic use of natural hot springs.
An consists mostly of outdoor pools (), which are sometimes at different temperatures. Extremely hot springs, where even experienced or frequent hot-spring bathers can only stay a few minutes, are called ('hell'). Many also have saunas, spa treatments and therapy centers. The same rules apply in public baths as in private baths, with bathers required to wash and clean themselves before entering the water. In general, the Japanese bathe naked in bathhouses; bathing suits are not permissible.
Art motif
Bathing scenes were already in the Middle Ages a popular subject of painters. Most of the subjects were women shown nude, but the interest was probably less to the bathing itself rather than to provide the context for representing the nude figure. From the Middle Ages, illustrated books of the time contained such bathing scenes. Biblical and mythological themes which featured bathing were depicted by numerous painters. Especially popular themes included Bathsheba in the bath, in which she is observed by King David, and Susanna in the sight of lecherous old men.
In the High Middle Ages, public baths were a popular subject of painting, with rather clear depictions of sexual advances, which probably were not based on actual observations. During the Renaissance and Baroque periods, the gods and nymphs of Greek mythology were depicted bathing in allegorical paintings by artists such as Titian and François Boucher, both of whom painted the goddess Diana bathing. Artists continued to paint Biblical characters bathing, and also sometimes depicted contemporary women bathing in the river, an example being Rembrandt's Woman Bathing.
In the 19th century, the use of the bathing scene reached its high point in classicism, realism and impressionism. Oriental themes and harem and turkish baths scenes became popular. These scenes were based on the artists' imagination, because access by men to Islamic women was not generally permitted. In the second half of the century, artists increasingly eschewed the pretexts of mythology and exoticism, and painted contemporary western women bathing. Edgar Degas, for example, painted over 100 paintings with a bathing theme. The subject of Bathers remained popular in avant-garde circles at the outset of the 20th century.
| Biology and health sciences | Health and fitness | null |
730393 | https://en.wikipedia.org/wiki/Boulevard | Boulevard | A boulevard is a type of broad avenue planted with rows of trees, or in parts of North America, any urban highway or wide road in a commercial district.
Boulevards were originally circumferential roads following the line of former city walls.
In North American usage, boulevards may be wide, multi-lane thoroughfares divided with only a central median.
Etymology
The word boulevard is borrowed from French. In France, it originally meant the flat surface of a rampart, and later a promenade taking the place of a demolished fortification. It is a borrowing from the Dutch word 'bulwark'.
Notable examples
Asia
Azerbaijan
Baku Boulevard
Cambodia
Norodom Boulevard
Monivong Boulevard
Sihanouk Boulevard
India
M G Road
Anna Salai
Indira Gandhi Sarani
Marine Drive
Krishnaraja Boulevard
Rajpath
Necklace Road
Mahatma Gandhi Road
Foreshore Road
Indonesia
Jalan Jenderal Sudirman
Jalan M.H. Thamrin
Jalan Jenderal Gatot Subroto
Jalan H.R. Rasuna Said
Jalan Gajah Mada/Jalan Hayam Wuruk
Jalan Prof. Dr. Satrio
Vietnam
Thang Long Boulevard
Iran
Keshavarz Boulevard
Philippines
Roxas Boulevard
Shaw Boulevard
España Boulevard
Quezon Boulevard
Aurora Boulevard
Osmeña Boulevard
Australia and Oceania
Australia
St Kilda Road, Melbourne
Royal Parade, Melbourne
Victoria Parade, Melbourne
Flemington Road, Melbourne
Mount Alexander Road, Melbourne
The Boulevard, Perth
New Zealand
The Four Avenues, Christchurch
Anzac Avenue, Dunedin
Marine Parade, Napier
Europe
Austria
Vienna Ring Road
Denmark
Boulevards in Copenhagen:
Nørre Voldgade
H. C. Andersens Boulevard
Vesterbrogade
Sønder Boulevard
Dalgas Boulevard
Strandboulevarden
Frederiksberg Allé
France
Boulevard Beaumarchais
Boulevard du Temple
Boulevard Montmartre
Boulevard des Italiens
Boulevard des Capucines
Boulevard de la Madeleine
Germany
Unter den Linden, Berlin
Kurfürstendamm, Berlin
Karl-Marx-Allee, Berlin
, Berlin
Königsallee, Düsseldorf
Goerdelerring, Leipzig
Brienner Straße, Munich
Leopoldstraße, Munich
Maximilianstraße, Munich
Prinzregentenstraße, Munich
Hungary
Little Boulevard, Budapest
Grand Boulevard, Budapest
Hungária Boulevard, Budapest
Ireland
O'Connell Street, Dublin
Italy
Viali di Circonvallazione, Florence
, Milan
Via Merulana, Rome
Netherlands
Lange Voorhout, The Hague
Spain
Gran Via de les Corts Catalanes, Barcelona
, Barcelona
Avinguda Diagonal, Barcelona
Alameda Principal, Málaga
Portugal
Avenida da Liberdade, Lisbon
Russia
Boulevard Ring, Moscow
Tverskoy Boulevard
Garden Ring, Moscow
North America
Canada
Lake Shore Boulevard, Toronto
Pie-IX Boulevard, Montreal
King George Boulevard, Surrey
Mexico
Paseo de la Reforma, Mexico City
United States
Ocean Parkway, Brooklyn
Broadway, Manhattan
West Side Highway, Manhattan
FDR Drive, Manhattan (future project)
Sunset Boulevard, Los Angeles
Santa Monica Boulevard, Los Angeles
Wilshire Boulevard, Los Angeles
Hollywood Boulevard, Los Angeles
Chicago Boulevard System
Benjamin Franklin Parkway, Philadelphia
Roosevelt Boulevard, Philadelphia
Southern Boulevard Park, Philadelphia
Boulevard, Atlanta
Park Avenue, New York City
Las Vegas Boulevard, Las Vegas
Ohio River Boulevard, Pittsburgh
South America
Argentina
, Buenos Aires
Uruguay
Artigas Boulevard, Montevideo
| Technology | Road infrastructure | null |
18298594 | https://en.wikipedia.org/wiki/Crossing%20number%20%28graph%20theory%29 | Crossing number (graph theory) | In graph theory, the crossing number of a graph is the lowest number of edge crossings of a plane drawing of the graph . For instance, a graph is planar if and only if its crossing number is zero. Determining the crossing number continues to be of great importance in graph drawing, as user studies have shown that drawing graphs with few crossings makes it easier for people to understand the drawing.
The study of crossing numbers originated in Turán's brick factory problem, in which Pál Turán asked for a factory plan that minimized the number of crossings between tracks connecting brick kilns to storage sites. Mathematically, this problem can be formalized as asking for the crossing number of a complete bipartite graph. The same problem arose independently in sociology at approximately the same time, in connection with the construction of sociograms. Turán's conjectured formula for the crossing numbers of complete bipartite graphs remains unproven, as does an analogous formula for the complete graphs.
The crossing number inequality states that, for graphs where the number of edges is sufficiently larger than the number of vertices, the crossing number is at least proportional to . It has applications in VLSI design and incidence geometry.
Without further qualification, the crossing number allows drawings in which the edges may be represented by arbitrary curves. A variation of this concept, the rectilinear crossing number, requires all edges to be straight line segments, and may differ from the crossing number. In particular, the rectilinear crossing number of a complete graph is essentially the same as the minimum number of convex quadrilaterals determined by a set of points in general position. The problem of determining this number is closely related to the happy ending problem.
Definitions
For the purposes of defining the crossing number, a drawing of an undirected graph is a mapping from the vertices of the graph to disjoint points in the plane, and from the edges of the graph to curves connecting their two endpoints. No vertex should be mapped onto an edge that it is not an endpoint of, and whenever two edges have curves that intersect (other than at a shared endpoint) their intersections should form a finite set of proper crossings, where the two curves are transverse. A crossing is counted separately for each of these crossing points, for each pair of edges that cross. The crossing number of a graph is then the minimum, over all such drawings, of the number of crossings in a drawing.
Some authors add more constraints to the definition of a drawing, for instance that each pair of edges have at most one intersection point (a shared endpoint or crossing). For the crossing number as defined above, these constraints make no difference, because a crossing-minimal drawing cannot have edges with multiple intersection points. If two edges with a shared endpoint cross, the drawing can be changed locally at the crossing point, leaving the rest of the drawing unchanged, to produce a different drawing with one fewer crossing. And similarly, if two edges cross two or more times, the drawing can be changed locally at two crossing points to make a different drawing with two fewer crossings. However, these constraints are relevant for variant definitions of the crossing number that, for instance, count only the numbers of pairs of edges that cross rather than the number of crossings.
Special cases
As of April 2015, crossing numbers are known for very few graph families. In particular, except for a few initial cases, the crossing number of complete graphs, bipartite complete graphs, and products of cycles all remain unknown, although there has been some progress on lower bounds.
Complete bipartite graphs
During World War II, Hungarian mathematician Pál Turán was forced to work in a brick factory, pushing wagon loads of bricks from kilns to storage sites. The factory had tracks from each kiln to each storage site, and the wagons were harder to push at the points where tracks crossed each other, from which Turán was led to ask his brick factory problem: how can the kilns, storage sites, and tracks be arranged to minimize the total number of crossings? Mathematically, the kilns and storage sites can be formalized as the vertices of a complete bipartite graph, with the tracks as its edges. A factory layout can be represented as a drawing of this graph, so the problem becomes:
what is the minimum possible number of crossings in a drawing of a complete bipartite graph?
Kazimierz Zarankiewicz attempted to solve Turán's brick factory problem; his proof contained an error, but he established a valid upper bound of
for the crossing number of the complete bipartite graph . This bound has been conjectured to be the optimal number of crossings for all complete bipartite graphs.
Complete graphs and graph coloring
The problem of determining the crossing number of the complete graph was first posed by Anthony Hill, and appeared in print in 1960. Hill and his collaborator John Ernest were two constructionist artists fascinated by mathematics. They not only formulated this problem but also originated a conjectural formula for this crossing number, which Richard K. Guy published in 1960. Namely, it is known that there always exists a drawing with
crossings. This formula gives values of for ; see sequence in the On-line Encyclopedia of Integer Sequences.
The conjecture is that there can be no better drawing, so that this formula gives the optimal number of crossings for the complete graphs. An independent formulation of the same conjecture was made by Thomas L. Saaty in 1964.
Saaty further verified that this formula gives the optimal number of crossings for and Pan and Richter showed that it also is optimal for .
The Albertson conjecture, formulated by Michael O. Albertson in 2007, states that, among all graphs with chromatic number , the complete graph has the minimum number of crossings. That is, if the conjectured formula for the crossing number of the complete graph is correct, then every -chromatic graph has crossing number at least equal to the same formula. The Albertson conjecture is now known to hold for .
Cubic graphs
The smallest cubic graphs with crossing numbers 1–11 are known . The smallest 1-crossing cubic graph is the complete bipartite graph , with 6 vertices. The smallest 2-crossing cubic graph is the Petersen graph, with 10 vertices. The smallest 3-crossing cubic graph is the Heawood graph, with 14 vertices. The smallest 4-crossing cubic graph is the Möbius-Kantor graph, with 16 vertices. The smallest 5-crossing cubic graph is the Pappus graph, with 18 vertices. The smallest 6-crossing cubic graph is the Desargues graph, with 20 vertices. None of the four 7-crossing cubic graphs, with 22 vertices, are well known. The smallest 8-crossing cubic graphs include the Nauru graph and the McGee graph or (3,7)-cage graph, with 24 vertices. The smallest 11-crossing cubic graphs include the Coxeter graph with 28 vertices.
In 2009, Pegg and Exoo conjectured that the smallest cubic graph with crossing number 13 is the Tutte–Coxeter graph and the smallest cubic graph with crossing number 170 is the Tutte 12-cage.
Connections to the bisection width
The 2/3-bisection width of a simple graph is the minimum number of edges whose removal results in a partition of the vertex set into two separated sets so that no set has more than vertices. Computing is NP-hard. Leighton proved that , provided that has bounded vertex degrees. This fundamental inequality can be used to derive an asymptotic lower bound for , when , or an estimate of it is known. In addition, this inequality has algorithmic application. Specifically, Bhat and Leighton used it (for the first time) for deriving an upper bound on the number of edge crossings in a drawing which is obtained by a divide and conquer approximation algorithm for computing .
Complexity and approximation
In general, determining the crossing number of a graph is hard; Garey and Johnson showed in 1983 that it is an NP-hard problem. In fact the problem remains NP-hard even when restricted to cubic graphs and to near-planar graphs (graphs that become planar after removal of a single edge). A closely related problem, determining the rectilinear crossing number, is complete for the existential theory of the reals.
On the positive side, there are efficient algorithms for determining whether the crossing number is less than a fixed constant . In other words, the problem is fixed-parameter tractable. It remains difficult for larger , such as . There are also efficient approximation algorithms for approximating on graphs of bounded degree which use the general and previously developed framework of Bhat and Leighton. In practice heuristic algorithms are used, such as the simple algorithm which starts with no edges and continually adds each new edge in a way that produces the fewest additional crossings possible. These algorithms are used in the Rectilinear Crossing Number distributed computing project.
The crossing number inequality
For an undirected simple graph with vertices and edges such that the crossing number is always at least
This relation between edges, vertices, and the crossing number was discovered independently by Ajtai, Chvátal, Newborn, and Szemerédi, and by Leighton . It is known as the crossing number inequality or crossing lemma.
The constant is the best known to date, and is due to Ackerman. The constant can be lowered to , but at the expense of replacing with the worse constant of .
The motivation of Leighton in studying crossing numbers was for applications to VLSI design in theoretical computer science. Later, Székely also realized that this inequality yielded very simple proofs of some important theorems in incidence geometry, such as Beck's theorem and the Szemerédi-Trotter theorem, and Tamal Dey used it to prove upper bounds on geometric k-sets.
Variations
If edges are required to be drawn as straight line segments, rather than arbitrary curves, then some graphs need more crossings. The rectilinear crossing number is defined to be the minimum number of crossings of a drawing of this type. It is always at least as large as the crossing number, and is larger for some graphs. It is known that, in general, the rectilinear crossing number can not be bounded by a function of the crossing number. The rectilinear crossing numbers for through are , () and values up to are known, with requiring either 7233 or 7234 crossings. Further values are collected by the Rectilinear Crossing Number project.
A graph has local crossing number if it can be drawn with at most crossings per edge, but not fewer.
The graphs that can be drawn with at most crossings per edge are also called -planar.
Other variants of the crossing number include the pairwise crossing number (the minimum number of pairs of edges that cross in any drawing) and the odd crossing number (the number of pairs of edges that cross an odd number of times in any drawing). The odd crossing number is at most equal to the pairwise crossing number, which is at most equal to the crossing number. However, by the Hanani–Tutte theorem, whenever one of these numbers is zero, they all are. surveys many such variants.
| Mathematics | Graph theory | null |
4590134 | https://en.wikipedia.org/wiki/Flathead%20grey%20mullet | Flathead grey mullet | The flathead grey mullet (Mugil cephalus) is an important food fish species in the mullet family Mugilidae. It is found in coastal temperate, tropical and subtropical waters worldwide. Its length is typically . It is known with numerous English names, including the flathead mullet, striped mullet (US, American Fisheries Society name), black mullet, bully mullet, common mullet, grey mullet, sea mullet and mullet, among others.
The flathead grey mullet is a mainly diurnal coastal species that often enters estuaries and rivers. It usually schools over sand or mud bottoms, feeding on zooplankton, dead plant matter, microalgae and detritus. The adult fish normally feed on algae in fresh water. The species is euryhaline, meaning that the fish can acclimate to different levels of salinity.
Description
The back of the fish is olive-green, sides are silvery and shade to white towards the belly. The fish may have six to seven distinctive lateral horizontal stripes. Lips are thin. The mullet has no lateral line. A common length is about , and its maximum length is . It can reach a maximum weight of .
Distribution
The flathead mullet is cosmopolitan in coastal waters of the tropical, subtropical and temperate zones of all seas, as far north as the Bay of Biscay and Nova Scotia in the Atlantic Ocean.
It occupies fresh, brackish and marine habitats in depths ranging between and with temperatures between .
In Australia, the fish is widespread, from Far North Queensland, around southern Australia to the Kimberley region of Western Australia. They also occur in the Bass Strait area of Tasmania. They live in tropical and temperate coastal marine and estuarine waters, but are also often found in the lower reaches of rivers. They are able to live in a wide range of salinity and so may also be found in lagoons, lakes and far into estuaries, but migrate back to the sea to spawn.
In freshwaters of the western United States, the striped mullet historically ranged far up the Colorado River to the vicinity of Blythe and up the Gila River to perhaps Tacna. Because of the dams and restricted flows to the Gulf of California, the range in Arizona is restricted to the Colorado River below Laguna Dam and the lower end of the Gila River when there is water present. They are often abundant in the mainstream and lateral canals in the Gila River region.
In the Colorado River mullet are pelagic in larger pools, sometimes moving into currents below dams, and generally occurring in small groups.
The mullet populations are currently declining in Arizona, due to periods when the Colorado River does not reach the Gulf of California.
Fisheries and aquaculture
The flathead grey mullet is an important food fish around the world, and it is both fished and farmed. The reported worldwide catches from fishing in 2012 were about 130,000 tonnes and aquaculture production was 142,000 tonnes.
Development
The ontogeny of mugilid larvae has been well studied, with the larval development of Mugil cephalus in particular being studied intensively due to its wide range of distribution and interest to aquaculture. The previously understudied osteological development of Mugil cephalus was investigated in a 2021 study, with four embryonic and six larval developmental steps being described in aquaculture-reared and wild-caught specimens. These descriptions provided clarification of questionable characters of adult mullets and revealed informative details with potential implications for phylogenetic hypotheses, as well as providing an overdue basis of comparison for aquaculture-reared mullets to enable recognition of malformations.
Cuisine
The roe of this mullet is salted, dried, and compressed to make a specialty food across the world, such as Greek avgotaraho, Taiwanese Wuyutsu, Korean eoran, Japanese karasumi, Italian bottarga, French poutargue, Turkish Haviar and Egyptian batarekh. In Egypt, the fish itself is salted, dried, and pickled to make fesikh.
On the coast of Northwest Florida and Alabama, this mullet, called the striped or black mullet, is often a specialty of seafood restaurants. Fried mullet is most popular, but smoked, baked, and canned mullet are also eaten. Local fishermen usually catch mullet in a castnet, though most use a land-based seine net. Mullet is a delicacy in this area and is most often consumed in the home. Mullet are usually filleted, and the remaining frames used for fish stock in chowders and stews. The mullet most commonly consumed in Florida however is the white mullet (Mugil curema), because its preference for cleaner water gives it a cleaner and less muddy taste.
| Biology and health sciences | Acanthomorpha | Animals |
4595284 | https://en.wikipedia.org/wiki/Viviparous%20brotula | Viviparous brotula | The viviparous brotulas form a family, the Bythitidae, of ophidiiform fishes. They are known as viviparous brotulas as they generally bear live young, although there are indications that some species (at least Didymothallus criniceps) do not. They are generally infrequently seen, somewhat tadpole-like in overall shape and mostly about in length, but some species grow far larger and may surpass .
Although many live near the coast in tropical or subtropical oceans, there are also species in deep water and cold oceans, for example Bythites. Thermichthys hollisi, which lives at depths of around , is associated with thermal vents. A few are fresh or brackish water cavefish: the Mexican blind brotula (Typhliasina pearsei), Galapagos cuskeel (Ogilbia galapagosensis), Diancistrus typhlops and some Lucifuga species.
Since 2002, more than 110 new species have been added to this family.
In 2005, 26 new species were described in a single paper by Danish and German scientists and in 2007, an additional eight new genera with 20 new species were described in another paper by the same scientists.
In some classifications the family Aphyonidae is placed within the Bythitidae and the tribe Dinematichthyini of the subfamily Brosmophycinae has been raised to the status of a family, the Dinematichthyidae which contains 25 genera and 114 species.
The Bythitidae is divided as follows:
Subfamily Brosmophycinae
Tribe Dinematichthyini
Alionematichthys
Beaglichthys
Brosmolus
Brotulinella
Dactylosurculus
Dermatopsis
Dermatopsoides
Diancistrus
Didymothallus
Dinematichthys
Dipulus
Gunterichthys
Lapitaichthys
Majungaichthys
Mascarenichthys
Monothrix
Nielsenichthys
Ogilbia
Ogilbichthys
Paradiancistrus
Porocephalichthys
Typhliasina
Ungusurculus
Zephyrichthys
Tribe Brosmophycini
Bidenichthys
Brosmodorsalis
Brosmophyciops
Brosmophycis
Eusurculus
Fiordichthys
Lucifuga
Melodichthys
Subfamily Bythitinae
Acarobythites
Anacanthobythites
Bellottia
Bythites
Calamopteryx
Cataetyx
Diplacanthopoma
Ematops
Grammonus
Hastatobythites
Hephthocara
Microbrotula
Parasaccogaster
Pseudogilbia
Pseudonus
Saccogaster
Stygnobrotula
Thalassobathia
Thermichthys
Timorichthys
Tuamotuichthys
| Biology and health sciences | Acanthomorpha | Animals |
7934681 | https://en.wikipedia.org/wiki/Gal%C3%A1pagos%20tortoise | Galápagos tortoise | The Galápagos tortoise or Galápagos giant tortoise (Chelonoidis niger) is a very large species of tortoise in the genus Chelonoidis (which also contains three smaller species from mainland South America). The species comprises 15 subspecies (13 extant and 2 extinct). It is the largest living species of tortoise, and can weigh up to . They are also the largest extant terrestrial cold-blooded animals (ectotherms).
With lifespans in the wild of over 100 years, it is one of the longest-lived vertebrates. Captive Galapagos tortoises can live up to 177 years. For example, a captive individual, Harriet, lived for at least 175 years. Spanish explorers, who discovered the islands in the 16th century, named them after the Spanish galápago, meaning "tortoise".
Galápagos tortoises are native to seven of the Galápagos Islands. Shell size and shape vary between subspecies and populations. On islands with humid highlands, the tortoises are larger, with domed shells and short necks; on islands with dry lowlands, the tortoises are smaller, with "saddleback" shells and long necks. Charles Darwin's observations of these differences on the second voyage of the Beagle in 1835, contributed to the development of his theory of evolution.
Tortoise numbers declined from over 250,000 in the 16th century to a low of around 15,000 in the 1970s. This decline was caused by overexploitation of the subspecies for meat and oil, habitat clearance for agriculture, and introduction of non-native animals to the islands, such as rats, goats, and pigs. The extinction of most giant tortoise lineages is thought to have also been caused by predation by humans or human ancestors, as the tortoises themselves have no natural predators. Tortoise populations on at least three islands have become extinct in historical times due to human activities. Specimens of these extinct taxa exist in several museums and also are being subjected to DNA analysis. 12 subspecies of the original 14–15 survive in the wild; a 13th subspecies (C. n. abingdonii) had only a single known living individual, kept in captivity and nicknamed Lonesome George until his death in June 2012. Two other subspecies, C. n. niger (the type subspecies of Galápagos tortoise) from Floreana Island and an undescribed subspecies from Santa Fe Island are known to have gone extinct in the mid-late 19th century. Conservation efforts, beginning in the 20th century, have resulted in thousands of captive-bred juveniles being released onto their ancestral home islands, and the total number of the subspecies is estimated to have exceeded 19,000 at the start of the 21st century. Despite this rebound, all surviving subspecies are classified as Threatened by the International Union for Conservation of Nature.
The Galápagos tortoises are one of two insular radiations of giant tortoises that still survive to the modern day; the other is Aldabrachelys gigantea of Aldabra and the Seychelles in the Indian Ocean, east of Tanzania. While giant tortoise radiations were common in prehistoric times, humans have wiped out the majority of them worldwide; the only other radiation of tortoises to survive to historic times, Cylindraspis of the Mascarenes, was driven to extinction by the 19th century, and other giant tortoise radiations such as a Centrochelys radiation on the Canary Islands and another Chelonoidis radiation in the Caribbean were driven to extinction prior to that.
Taxonomy
Early classification
The Galápagos Islands were discovered in 1535, but first appeared on the maps, of Gerardus Mercator and Abraham Ortelius, around 1570. The islands were named "Insulae de los Galopegos" (Islands of the Tortoises) in reference to the giant tortoises found there.
Initially, the giant tortoises of the Indian Ocean and those from the Galápagos were thought to be the same subspecies. Naturalists thought that sailors had transported the tortoises there. In 1676, the pre-Linnaean authority Claude Perrault referred to both subspecies as Tortue des Indes. In 1783, Johann Gottlob Schneider classified all giant tortoises as Testudo indica ("Indian tortoise"). In 1812, August Friedrich Schweigger named them Testudo gigantea ("gigantic tortoise"). In 1834, André Marie Constant Duméril and Gabriel Bibron classified the Galápagos tortoises as a separate subspecies, which they named Testudo nigrita ("black tortoise").
Recognition of subpopulations
The first systematic survey of giant tortoises was by the zoologist Albert Günther of the British Museum, in 1875. Günther identified at least five distinct populations from the Galápagos, and three from the Indian Ocean islands. He expanded the list in 1877 to six from the Galápagos, four from the Seychelles, and four from the Mascarenes. Günther hypothesized that all the giant tortoises descended from a single ancestral population which had spread by sunken land bridges. This hypothesis was later disproven by the understanding that the Galápagos, the atolls of Seychelles, and the Mascarene islands are all of recent volcanic origin and have never been linked to a continent by land bridges. Galápagos tortoises are now thought to have descended from a South American ancestor, while the Indian Ocean tortoises derived from ancestral populations on Madagascar.
At the end of the 19th century, Georg Baur and Walter Rothschild recognised five more populations of Galápagos tortoise. In 1905–06, an expedition by the California Academy of Sciences, with Joseph R. Slevin in charge of reptiles, collected specimens which were studied by Academy herpetologist John Van Denburgh. He identified four additional populations, and proposed the existence of 15 subspecies. Van Denburgh's list still guides the taxonomy of the Galápagos tortoise, though now 10 populations are thought to have existed.
Current species and genus names
The current specific designation of niger, formerly feminized to nigra ("black"—Quoy & Gaimard, 1824b) was resurrected in 1984 after it was discovered to be the senior synonym (an older taxonomic synonym taking historical precedence) for the then commonly used subspecies name of elephantopus ("elephant-footed" – Harlan, 1827). Quoy and Gaimard's Latin description explains the use of nigra: "Testudo toto corpore nigro" means "tortoise with completely black body". Quoy and Gairmard described nigra from a living specimen, but no evidence indicates they knew of its accurate provenance within the Galápagos – the locality was in fact given as California. Garman proposed the linking of nigra with the extinct Floreana subspecies. Later, Pritchard deemed it convenient to accept this designation, despite its tenuousness, for minimal disruption to the already confused nomenclature of the subspecies. The even more senior subspecies synonym of californiana ("californian" – Quoy & Gaimard, 1824a) is considered a nomen oblitum ("forgotten name").
Previously, the Galápagos tortoise was considered to belong to the genus Geochelone, known as 'typical tortoises' or 'terrestrial turtles'. In the 1990s, subgenus Chelonoidis was elevated to generic status based on phylogenetic evidence which grouped the South American members of Geochelone into an independent clade (branch of the tree of life). This nomenclature has been adopted by several authorities.
Subspecies
Within the archipelago, 14-15 subspecies of Galápagos tortoises have been identified, although only 12 survive to this day. Five are found on separate islands; five of them on the volcanoes of Isabela Island. Several of the surviving subspecies are seriously endangered. A 13th subspecies, C. n. abingdonii from Pinta Island, is extinct since 2012. The last known specimen, named Lonesome George, died in captivity on 24 June 2012; George had been mated with female tortoises of several other subspecies, but none of the eggs from these pairings hatched. The subspecies inhabiting Floreana Island (C. niger) is thought to have been hunted to extinction by 1850, only 15 years after Charles Darwin's landmark visit of 1835, when he saw shells, but no live tortoises there. However, recent DNA testing shows that an intermixed, non-native population currently existing on the island of Isabela is of genetic resemblance to the subspecies native to Floreana, suggesting that C. niger has not gone entirely extinct. The existence of the C. n. phantastica subspecies of Fernandina Island was disputed, as it was described from a single specimen that may have been an artificial introduction to the island; however, a live female was found in 2019, likely confirming the subspecies' validity.
Prior to widespread knowledge of the differences between the populations (sometimes called races) from different islands and volcanoes, captive collections in zoos were indiscriminately mixed. Fertile offspring resulted from pairings of animals from different races. However, captive crosses between tortoises from different races have lower fertility and higher mortality than those between tortoises of the same race, and captives in mixed herds normally direct courtship only toward members of the same race.
The valid scientific names of each of the individual populations are not universally accepted, and some researchers still consider each subspecies to be distinct species. Prior to 2021, all subspecies were classified as distinct species from one another, but a 2021 study analyzing the level of divergence within the extinct West Indian Chelonoidis radiation and comparing it to the Galápagos radiation found that the level of divergence within both clades may have been significantly overestimated, and supported once again reclassifying all Galápagos tortoises as subspecies of a single species, C. niger. This was followed by the Turtle Taxonomy Working Group and the Reptile Database later that year. The taxonomic status of the various races is not fully resolved.
Testudo californiana Quoy & Gaimard, 1824a (nomen oblitum)
Testudo nigra Quoy & Gaimard, 1824b (nomen novum)
Testudo elephantopus Harlan, 1827 (nomen dubium)
Testudo nigrita Duméril and Bibron, 1834 (nomen dubium)
Testudo planiceps Gray, 1853 (nomen dubium)
Testudo clivosa Garman, 1917 (nomen dubium)
Testudo typica Garman, 1917 (nomen dubium)
Testudo (Chelonoidis) elephantopus Williams, 1952
Geochelone (Chelonoidis) elephantopus Pritchard, 1967
Chelonoidis elephantopus Bour, 1980
C. n. nigra (nominate subspecies)
Testudo californiana Quoy & Gaimard, 1824a (nomen oblitum)
Testudo nigra Quoy & Gaimard, 1824b (nomen novum)
Testudo galapagoensis Baur 1889
C. n. abingdoni
Testudo ephippiumGünther, 1875 (partim, misidentified type specimen once erroneously attributed to what is now C. n. duncanensis)
Testudo abingdoniGünther, 1877
C. n. becki
Testudo beckiRothschild, 1901
C. n. chathamensis
Testudo wallaceiRothschild 1902 (partim, nomen dubium)
Testudo chathamensisVan Denburgh, 1907
C. n. darwini
Testudo wallaceiRothschild 1902 (partim, nomen dubium)
Testudo darwiniVan Denburgh, 1907
C. n. duncanensis
Testudo ephippiumGünther, 1875 (partim, misidentified type)
Geochelone nigra duncanensisGarman, 1917 in Pritchard, 1996(nomen nudum)
C. n. hoodensis
Testudo hoodensisVan Denburgh, 1907
C. n. phantastica
Testudo phantasticusVan Denburgh, 1907
C. n. porteri
Testudo nigrita Duméril and Bibron, 1834 (nomen dubium)
Testudo porteriRothschild, 1903
C. n. vicina
Testudo microphyesGünther, 1875
Testudo vicinaGünther, 1875
Testudo güntheriBaur, 1889
Testudo macrophyesGarman, 1917
Testudo vandenburghiDe Sola, R. 1930 (nomen nudum)
Chelonoidis nigra nigra
Testudo nigra Quoy & Gaimard, 1824
Testudo californiana Quoy & Gaimard, 1824
Testudo galapagoensis Baur, 1889
Testudo elephantopus galapagoensis Mertens & Wermuth, 1955
Geochelone elephantopus galapagoensis Pritchard, 1967
Chelonoidis galapagoensis Bour, 1980
Chelonoidis nigra Bour, 1985
Chelonoidis elephantopus galapagoensis Obst, 1985
Geochelone nigra Pritchard, 1986
Geochelone nigra nigra Stubbs, 1989
Chelonoidis nigra galapagoensis David, 1994
Chelonoidis nigra nigra David, 1994
Geochelone elephantopus nigra Bonin, Devaux & Dupré, 1996
Testudo california Paull, 1998 (ex errore)
Testudo californianana Paull, 1999 (ex errore)
Chelonoidis nigra abingdonii
Testudo ephippium Günther, 1875
Testudo abingdonii Günther, 1877
Testudo abingdoni Van Denburgh, 1914 (ex errore)
Testudo elephantopus abingdonii Mertens & Wermuth, 1955
Testudo elephantopus ephippium Mertens & Wermuth, 1955
Geochelone abingdonii Pritchard, 1967
Geochelone elephantopus abingdoni Pritchard, 1967
Geochelone elephantopus ephippium Pritchard, 1967
Geochelone ephippium Pritchard, 1967
Chelonoidis abingdonii Bour, 1980
Chelonoidis ephippium Bour, 1980
Geochelone elephantopus abingdonii Groombridge, 1982
Geochelone abingdoni Fritts, 1983
Geochelone epphipium Fritts, 1983 (ex errore)
Chelonoidis nigra ephippium Pritchard, 1984
Chelonoidis elephantopus abingdoni Obst, 1985
Chelonoidis elephantopus ephippium Obst, 1985
Geochelone nigra abingdoni Stubbs, 1989
Chelonoidis nigra abingdonii David, 1994
Chelonoidis elephantopus abingdonii Rogner, 1996
Chelonoidis nigra abingdonii Bonin, Devaux & Dupré, 1996
Chelonoidis nigra abdingdonii Obst, 1996 (ex errore)
Geochelone abdingdonii Obst, 1996
Geochelone nigra abdingdoni Obst, 1996 (ex errore)
Geochelone nigra ephyppium Caccone, Gibbs, Ketmaier, Suatoni & Powell, 1999 (ex errore)
Chelonoidis nigra ahingdonii Artner, 2003 (ex errore)
Chelonoidis abingdoni Joseph-Ouni, 2004
Chelonoidis nigra becki
Testudo becki Rothschild, 1901
Testudo bedsi Heller, 1903 (ex errore)
Geochelone becki Pritchard, 1967
Geochelone elephantopus becki Pritchard, 1967
Chelonoidis becki Bour, 1980
Chelonoidis elephantopus becki Obst, 1985
Chelonoidis nigra beckii David, 1994 (ex errore)
Chelonoidis elephantopus beckii Rogner, 1996
Chelonoidis nigra becki Obst, 1996
Chelonoidis nigra chathamensis
Testudo wallacei Rothschild, 1902
Testudo chathamensis Van Denburgh, 1907
Testudo elephantopus chathamensis Mertens & Wermuth, 1955
Testudo elephantopus wallacei Mertens & Wermuth, 1955
Testudo chatamensis Slevin & Leviton, 1956 (ex errore)
Geochelone chathamensis Pritchard, 1967
Geochelone elephantopus chathamensis Pritchard, 1967
Geochelone elephantopus wallacei Pritchard, 1967
Geochelone wallacei Pritchard, 1967
Chelonoidis chathamensis Bour, 1980
Chelonoidis elephantopus chathamensis Obst, 1985
Chelonoidis elephantopus wallacei Obst, 1985
Chelonoidis elephantopus chatamensis Gosławski & Hryniewicz, 1993
Chelonoidis nigra chathamensis David, 1994
Chelonoidis nigra wallacei Bonin, Devaux & Dupré, 1996
Geochelone cathamensis Obst, 1996 (ex errore)
Geochelone elephantopus chatamensis Paull, 1996
Testudo chathamensis chathamensis Pritchard, 1998
Cherlonoidis nigra wallacei Wilms, 1999
Geochelone nigra chatamensis Caccone, Gibbs, Ketmaier, Suatoni & Powell, 1999
Geochelone nigra wallacei Chambers, 2004
Chelonoidis nigra darwini
Testudo wallacei Rothschild, 1902
Testudo darwini Van Denburgh, 1907
Testudo elephantopus darwini Mertens & Wermuth, 1955
Testudo elephantopus wallacei Mertens & Wermuth, 1955
Geochelone darwini Pritchard, 1967
Geochelone elephantopus darwini Pritchard, 1967
Geochelone elephantopus wallacei Pritchard, 1967
Geochelone wallacei Pritchard, 1967
Chelonoidis darwini Bour, 1980
Chelonoidis elephantopus darwini Obst, 1985
Chelonoidis elephantopus wallacei Obst, 1985
Chelonoidis nigra darwinii David, 1994 (ex errore)
Chelonoidis elephantopus darwinii Rogner, 1996
Chelonoidis nigra darwini Bonin, Devaux & Dupré, 1996
Chelonoidis nigra wallacei Bonin, Devaux & Dupré, 1996
Cherlonoidis nigra wallacei Wilms, 1999
Geochelone nigra darwinii Ferri, 2002
Geochelone nigra wallacei Chambers, 2004
Chelonoidis nigra duncanensis
Testudo duncanensis Garman, 1917 (nomen nudum)
Geochelone nigra duncanensis Stubbs, 1989
Geochelone nigra duncanensis Garman, 1996
Chelonoidis nigra duncanensis Artner, 2003
Chelonoidis duncanensis Joseph-Ouni, 2004
Chelonoidis nigra hoodensis
Testudo hoodensis Van Denburgh, 1907
Testudo elephantopus hoodensis Mertens & Wermuth, 1955
Geochelone elephantopus hoodensis Pritchard, 1967
Geochelone hoodensis Pritchard, 1967
Chelonoidis hoodensis Bour, 1980
Chelonoidis elephantopus hoodensis Obst, 1985
Chelonoidis nigra hoodensis David, 1994
Chelonoidis nigra phantastica
Testudo phantasticus Van Denburgh, 1907
Testudo phantastica Siebenrock, 1909
Testudo elephantopus phantastica Mertens & Wermuth, 1955
Geochelone elephantopus phantastica Pritchard, 1967
Geochelone phantastica Pritchard, 1967
Chelonoidis phantastica Bour, 1980
Geochelone phantasticus Crumly, 1984
Chelonoidis elephantopus phantastica Obst, 1985
Chelonoidis nigra phantastica David, 1994
Chelonoidis nigra porteri
Testudo nigrita Duméril & Bibron, 1835
Testudo porteri Rothschild, 1903
Testudo elephantopus nigrita Mertens & Wermuth, 1955
Geochelone elephantopus porteri Pritchard, 1967
Geochelone nigrita Pritchard, 1967
Chelonoidis nigrita Bour, 1980
Geochelone elephantopus nigrita Honegger, 1980
Geochelone porteri Fritts, 1983
Chelonoidis elephantopus nigrita Obst, 1985
Geochelone nigra porteri Stubbs, 1989
Chelonoidis elephantopus porteri Gosławski & Hryniewicz, 1993
Chelonoidis nigra nigrita David, 1994
Geochelone nigra perteri Müller & Schmidt, 1995 (ex errore)
Chelonoidis nigra porteri Bonin, Devaux & Dupré, 1996
Chelonoidis nigra vicina
Testudo elephantopus Harlan, 1827
Testudo microphyes Günther, 1875
Testudo vicina Günther, 1875
Testudo macrophyes Garman, 1917
Testudo vandenburghi de Sola, 1930
Testudo elephantopus elephantopus Mertens & Wermuth, 1955
Geochelone elephantopus Williams, 1960
Geochelone elephantopus elephantopus Pritchard, 1967
Geochelone elephantopus guentheri Pritchard, 1967
Geochelone elephantopus guntheri Pritchard, 1967 (ex errore)
Geochelone elephantopus microphyes Pritchard, 1967
Geochelone elephantopus vandenburgi Pritchard, 1967 (ex errore)
Geochelone guntheri Pritchard, 1967
Geochelone microphyes Pritchard, 1967
Geochelone vandenburghi Pritchard, 1967
Geochelone vicina Pritchard, 1967
Geochelone elephantopus microphys Arnold, 1979 (ex errore)
Geochelone elephantopus vandenburghi Pritchard, 1979
Chelonoides elephantopus Obst, 1980
Chelonoidis elephantopus Bour, 1980
Chelonoidis guentheri Bour, 1980
Chelonoidis microphyes Bour, 1980
Chelonoidis vandenburghi Bour, 1980
Geochelone guentheri Fritts, 1983
Chelonoidis elephantopus elephantopus Obst, 1985
Chelonoidis elephantopus guentheri Obst, 1985
Chelonoidis elephantopus microphyes Obst, 1985
Chelonoidis elephantopus vandenburghi Obst, 1985
Geochelone elephantopus vicina Swingland, 1989
Geochelone elephantopus vicini Swingland, 1989 (ex errore)
Chelonoidis elephantopus guntheri Gosławski & Hryniewicz, 1993
Chelonoidis nigra guentheri David, 1994
Chelonoidis nigra microphyes David, 1994
Chelonoidis nigra vandenburghi David, 1994
Geochelone nigra elephantopus Müller & Schmidt, 1995
Chelonoidis elephantopus vicina Rogner, 1996
Geochelone elephantopus vandenburghii Obst, 1996 (ex errore)
Geochelone vandenburghii Obst, 1996
Chelonoidis nigra microphyies Bonin, Devaux & Dupré, 1996 (ex errore)
Geochelone elephantopus microphytes Paull, 1996 (ex errore)
Geochelone elephantopus vandenbergi Paull, 1996 (ex errore)
Testudo elephantopus guntheri Paull, 1999
Chelonoidis nigra vicina Artner, 2003
Chelonoidis vicina Joseph-Ouni, 2004
Geochelone nigra guentheri Chambers, 2004
Modern DNA methods have revealed new information on the relationships between the subspecies:
Isabela Island
The five populations living on the largest island, Isabela, are the ones that are the subject of the most debate as to whether they are true subspecies or just distinct populations of a subspecies. It is widely accepted that the population living on the northernmost volcano, Volcan Wolf, is genetically independent from the four populations to the south and is therefore a separate subspecies. It is thought to be derived from a different colonization event than the others. A colonization from the island of Santiago apparently gave rise to the Volcan Wolf subspecies (C. n. becki) while the four southern populations are believed to be descended from a second colonization from the more southerly island of Santa Cruz. Tortoises from Santa Cruz are thought to have first colonized the Sierra Negra volcano, which was the first of the island's volcanoes to form. The tortoises then spread north to each newly created volcano, resulting in the populations living on Volcan Alcedo and then Volcan Darwin. Recent genetic evidence shows that these two populations are genetically distinct from each other and from the population living on Sierra Negra (C. guentheri) and therefore form the subspecies C. n. vandenburghi (Alcedo) and C. n. microphyes (Darwin). The fifth population living on the southernmost volcano (C. n. vicina) is thought to have split off from the Sierra Negra population more recently and is therefore not as genetically different as the other two. Isabela is the most recently formed island tortoises inhabit, so its populations have had less time to evolve independently than populations on other islands, but according to some researchers, they are all genetically different and should each be considered as separate subspecies.
Floreana Island
Phylogenetic analysis may help to "resurrect" the extinct subspecies of Floreana (C. n. niger)—a subspecies known only from subfossil remains. Some tortoises from Isabela were found to be a partial match for the genetic profile of Floreana specimens from museum collections, possibly indicating the presence of hybrids from a population transported by humans from Floreana to Isabela, resulting either from individuals deliberately transported between the islands, or from individuals thrown overboard from ships to lighten the load. Nine Floreana descendants have been identified in the captive population of the Fausto Llerena Breeding Center on Santa Cruz; the genetic footprint was identified in the genomes of hybrid offspring. This allows the possibility of re-establishing a reconstructed subspecies from selective breeding of the hybrid animals. Furthermore, individuals from the subspecies possibly are still extant. Genetic analysis from a sample of tortoises from Volcan Wolf found 84 first-generation C. n. niger hybrids, some less than 15 years old. The genetic diversity of these individuals is estimated to have required 38 C. n. niger parents, many of which could still be alive on Isabela Island.
Pinta Island
The Pinta Island subspecies (C. n. abingdonii, now extinct) has been found to be most closely related to the subspecies on the islands of San Cristóbal (C. n. chathamensis) and Española (C. n. hoodensis) which lie over 300 km (190 mi) away, rather than that on the neighbouring island of Isabela as previously assumed. This relationship is attributable to dispersal by the strong local current from San Cristóbal towards Pinta. This discovery informed further attempts to preserve the C. n. abingdonii lineage and the search for an appropriate mate for Lonesome George, which had been penned with females from Isabela. Hope was bolstered by the discovery of a C. n. abingdonii hybrid male in the Volcán Wolf population on northern Isabela, raising the possibility that more undiscovered living Pinta descendants exist.
Santa Cruz Island
Mitochondrial DNA studies of tortoises on Santa Cruz show up to three genetically distinct lineages found in nonoverlapping population distributions around the regions of Cerro Montura, Cerro Fatal, and La Caseta. Although traditionally grouped into a single subspecies (C. n. porteri), the lineages are all more closely related to tortoises on other islands than to each other: Cerro Montura tortoises are most closely related to C. n. duncanensis from Pinzón, Cerro Fatal to C. n. chathamensis from San Cristóbal, and La Caseta to the four southern races of Isabela as well as Floreana tortoises.
In 2015, the Cerro Fatal tortoises were described as a distinct taxon, donfaustoi. Prior to the identification of this subspecies through genetic analysis, it was noted that there existed differences in shells between the Cerro Fatal tortoises and other tortoises on Santa Cruz. By classifying the Cerro Fatal tortoises into a new taxon, greater attention can be paid to protecting its habitat, according to Adalgisa Caccone, who is a member of the team making this classification.
Pinzón Island
When it was discovered that the central, small island of Pinzón had only 100–200 very old adults and no young tortoises had survived into adulthood for perhaps more than 70 years, the resident scientists initiated what would eventually become the Giant Tortoise Breeding and Rearing Program. Over the next 50 years, this program resulted in major successes in the recovery of giant tortoise populations throughout the archipelago.
In 1965, the first tortoise eggs collected from natural nests on Pinzón Island were brought to the Charles Darwin Research Station, where they would complete the period of incubation and then hatch, becoming the first young tortoises to be reared in captivity. The introduction of black rats onto Pinzón sometime in the latter half of the 19th century had resulted in the complete eradication of all young tortoises. Black rats had been eating both tortoise eggs and hatchlings, effectively destroying the future of the tortoise population. Only the longevity of giant tortoises allowed them to survive until the Galápagos National Park, Island Conservation, Charles Darwin Foundation, the Raptor Center, and Bell Laboratories removed invasive rats in 2012. In 2013, heralding an important step in Pinzón tortoise recovery, hatchlings emerged from native Pinzón tortoise nests on the island and the Galápagos National Park successfully returned 118 hatchlings to their native island home. Partners returned to Pinzón Island in late 2014 and continued to observe hatchling tortoises (now older), indicating that natural recruitment is occurring on the island unimpeded. They also discovered a snail subspecies new to science. These exciting results highlight the conservation value of this important management action. In early 2015, after extensive monitoring, partners confirmed that Pinzón and Plaza Sur Islands are now both rodent-free.
Española Island
On the southern island of Española, only 14 adult tortoises were found, two males and 12 females. The tortoises apparently were not encountering one another, so no reproduction was occurring. Between 1963 and 1974, all 14 adult tortoises discovered on the island were brought to the tortoise center on Santa Cruz and a tortoise breeding program was initiated. In 1977, a third Española male tortoise was returned to Galapagos from the San Diego Zoo and joined the breeding group. After 40 years' work reintroducing captive animals, a detailed study of the island's ecosystem has confirmed it has a stable, breeding population. Where once 15 were known, now more than 1,000 giant tortoises inhabit the island of Española. One research team has found that more than half the tortoises released since the first reintroductions are still alive, and they are breeding well enough for the population to progress onward, unaided. In January 2020, it was widely reported that Diego, a 100-year-old male tortoise, resurrected 40% of the tortoise population on the island and is known as the "Playboy Tortoise".
Fernandina Island
The C. n. phantasticus subspecies from Fernandina was originally known from a single specimen—an old male from the voyage of 1905–06. No other tortoises or remains were found on the island for a long time after its sighting, leading to suggestions that the specimen was an artificial introduction from elsewhere. Fernandina has neither human settlements nor feral mammals, so if this subspecies ever did exist, its extinction would have been by natural means, such as volcanic activity. Nevertheless, there have occasionally been reports from Fernandina. In 2019, an elderly female specimen was finally discovered on Fernandina and transferred to a breeding center, and trace evidence found on the expedition indicates that more individuals likely exist in the wild. It has been theorized that the rarity of the subspecies may be due to the harsh habitat it survives in, such as the lava flows that are known to frequently cover the island.
Santa Fe Island
The extinct Santa Fe subspecies has not yet been described and thus has no binomial name, having been identified from the limited evidence of bone fragments (but no shells, the most durable part) of 14 individuals, old eggs, and old dung found on the island in 1905–06. The island has never been inhabited by man nor had any introduced predators, but reports have been made of whalers hauling tortoises off the island. Later genetic studies of the bone fragments indicate that the Santa Fe subspecies was distinct, and was most closely related to C. n. hoodensis. A population of C. n. hoodensis has since been reintroduced to and established on the island to fill in the ecological role of the Santa Fe tortoise.
Subspecies of doubtful existence
The purported Rábida Island subspecies (C. n. wallacei) was described from a single specimen collected by the California Academy of Sciences in December 1905, which has since been lost. This individual was probably an artificial introduction from another island that was originally penned on Rábida next to a good anchorage, as no contemporary whaling or sealing logs mention removing tortoises from this island.
Description
The tortoises have a large bony shell of a dull brown or grey color. The plates of the shell are fused with the ribs in a rigid protective structure that is integral to the skeleton. Lichens can grow on the shells of these slow-moving animals. Tortoises keep a characteristic scute (shell segment) pattern on their shells throughout life, though the annual growth bands are not useful for determining age because the outer layers are worn off with time. A tortoise can withdraw its head, neck, and fore limbs into its shell for protection. The legs are large and stumpy, with dry, scaly skin and hard scales. The front legs have five claws, the back legs four.
Gigantism
The discoverer of the Galápagos Islands, Fray Tomás de Berlanga, Bishop of Panama, wrote in 1535 of "such big tortoises that each could carry a man on top of himself." Naturalist Charles Darwin remarked after his trip three centuries later in 1835, "These animals grow to an immense size ... several so large that it required six or eight men to lift them from the ground". The largest recorded individuals have reached weights of over and lengths of . Size overlap is extensive with the Aldabra giant tortoise, however taken as a subspecies, the Galápagos tortoise seems to average slightly larger, with weights in excess of being slightly more commonplace. Weights in the larger bodied subspecies range from in mature males and from in adult females. However, the size is variable across the islands and subspecies; those from Pinzón Island are relatively small with a maximum known weight of and carapace length of approximately compared to range in tortoises from Santa Cruz Island. The tortoises' gigantism was probably a trait useful on continents that was fortuitously helpful for successful colonisation of these remote oceanic islands rather than an example of evolved insular gigantism. Large tortoises would have a greater chance of surviving the journey over water from the mainland as they can hold their heads a greater height above the water level and have a smaller surface area/volume ratio, which reduces osmotic water loss. Their significant water and fat reserves would allow the tortoises to survive long ocean crossings without food or fresh water, and to endure the drought-prone climate of the islands. A larger size allowed them to better tolerate extremes of temperature due to gigantothermy. Fossil giant tortoises from mainland South America have been described that support this hypothesis of gigantism that pre-existed the colonization of islands.
Shell shape
Galápagos tortoises possess two main shell forms that correlate with the biogeographic history of the subspecies group. They exhibit a spectrum of carapace morphology ranging from "saddleback" (denoting upward arching of the front edge of the shell resembling a saddle) to "domed" (denoting a rounded convex surface resembling a dome). When a saddleback tortoise withdraws its head and forelimbs into its shell, a large unprotected gap remains over the neck, evidence of the lack of predation during the evolution of this structure. Larger islands with humid highlands over in elevation, such as Santa Cruz, have abundant vegetation near the ground. Tortoises native to these environments tend to have domed shells and are larger, with shorter necks and limbs. Saddleback tortoises originate from small islands less than in elevation with dry habitats (e.g. Española and Pinzón) that are more limited in food and other resources. Two lineages of Galápagos tortoises possess the Island of Santa Cruz and when observed it is concluded that despite the shared similarities of growth patterns and morphological changes observed during growth, the two lineages and two sexes can be distinguished on the basis of distinct carapace features. Lineages differ by the shape of the vertebral and pleural scutes. Females have a more elongated and wider carapace shape than males. Carapace shape changes with growth, with vertebral scutes becoming narrower and pleural scutes becoming larger during late ontogeny.
Evolutionary implications
In combination with proportionally longer necks and limbs, the unusual saddleback carapace structure is thought to be an adaptation to increase vertical reach, which enables the tortoise to browse tall vegetation such as the Opuntia (prickly pear) cactus that grows in arid environments. Saddlebacks are more territorial and smaller than domed varieties, possibly adaptations to limited resources. Alternatively, larger tortoises may be better-suited to high elevations because they can resist the cooler temperatures that occur with cloud cover or fog.
A competing hypothesis is that, rather than being principally a feeding adaptation, the distinctive saddle shape and longer extremities might have been a secondary sexual characteristic of saddleback males. Male competition over mates is settled by dominance displays on the basis of vertical neck height rather than body size (see below). This correlates with the observation that saddleback males are more aggressive than domed males. The shell distortion and elongation of the limbs and neck in saddlebacks is probably an evolutionary compromise between the need for a small body size in dry conditions and a high vertical reach for dominance displays.
The saddleback carapace probably evolved independently several times in dry habitats, since genetic similarity between populations does not correspond to carapace shape. Saddleback tortoises are, therefore, not necessarily more closely related to each other than to their domed counterparts, as shape is not determined by a similar genetic background, but by a similar ecological one.
Sexual dimorphism
Sexual dimorphism is most pronounced in saddleback populations in which males have more angled and higher front openings, giving a more extreme saddled appearance. Males of all varieties generally have longer tails and shorter, concave plastrons with thickened knobs at the back edge to facilitate mating. Males are larger than females—adult males weigh around while females are .
Behavior
Routine
The tortoises are ectothermic (cold-blooded), so they bask for 1–2 hours after dawn to absorb the sun's heat through their dark shells before actively foraging for 8–9 hours a day. They travel mostly in the early morning or late afternoon between resting and grazing areas. They have been observed to walk at a speed of .
On the larger and more humid islands, the tortoises seasonally migrate between low elevations, which become grassy plains in the wet season, and meadowed areas of higher elevation (up to ) in the dry season. The same routes have been used for many generations, creating well-defined paths through the undergrowth known as "tortoise highways". On these wetter islands, the domed tortoises are gregarious and often found in large herds, in contrast to the more solitary and territorial disposition of the saddleback tortoises.
Tortoises sometimes rest in mud wallows or rain-formed pools, which may be both a thermoregulatory response during cool nights, and a protection from parasites such as mosquitoes and ticks. Parasites are countered by taking dust baths in loose soil. Some tortoises have been noted to shelter at night under overhanging rocks. Others have been observed sleeping in a snug depression in the earth or brush called a "pallet". Local tortoises using the same pallet sites, such as on Volcán Alcedo, results in the formation of small, sandy pits.
Diet
The tortoises are herbivores that consume a diet of cacti, grasses, leaves, lichens, berries, melons, oranges and milkweed. They have been documented feeding on Hippomane mancinella (poison apple), the endemic guava Psidium galapageium, the water fern Azolla microphylla, bromeliad Tillandsia insularis and the Galápagos tomato Solanum cheesmaniae. Juvenile tortoises eat an average of 16.7% of their own body weight in dry matter per day, with a digestive efficiency roughly equal to that of hindgut-fermenting herbivorous mammals such as horses and rhinos.
Tortoises acquire most of their water from dew and sap in vegetation (particularly the Opuntia cactus), which enables them to survive for more than six months without drinking. They can endure up to a year when deprived of all food and water, surviving by breaking down their body fat to produce water as a byproduct. Tortoises also have very slow metabolisms. When thirsty, they may drink large quantities of water very quickly, storing it in their bladders and the "root of the neck" (the pericardium), both of which served to make them useful water sources on ships. On arid islands, tortoises lick morning dew from boulders, and the repeated action over many generations has formed half-sphere depressions in the rock.
Senses
Regarding their senses, Charles Darwin observed, "The inhabitants believe that these animals are absolutely deaf; certainly they do not overhear a person walking near behind them. I was always amused, when overtaking one of these great monsters as it was quietly pacing along, to see how suddenly, the instant I passed, it would draw in its head and legs, and uttering a deep hiss fall to the ground with a heavy sound, as if struck dead." Although they are not deaf, tortoises depend far more on vision and smell as stimuli than hearing.
Mutualism
Tortoises share a mutualistic relationship with some subspecies of Galápagos finch and mockingbirds. The birds benefit from the food source and the tortoises get rid of irritating ectoparasites.
Small groups of finches initiate the process by hopping on the ground in an exaggerated fashion facing the tortoise. The tortoise signals it is ready by rising up and extending its neck and legs, enabling the birds to reach otherwise inaccessible spots on the tortoise's body such as the neck, rear legs, cloacal opening, and skin between plastron and carapace.
Some tortoises have been observed to exploit this mutualistic relationship to consume birds seeking to groom them. After rising and extending its limbs, the bird may go beneath the tortoise to investigate, whereupon suddenly the tortoise withdraws its limbs to drop flat and kill the bird. It then steps back to eat the bird, presumably to supplement its diet with protein.
Mating
Mating occurs at any time of the year, although it does have seasonal peaks between February and June in the humid uplands during the rainy season. When mature males meet in the mating season, they face each other in a ritualised dominance display, rise up on their legs, and stretch up their necks with their mouths gaping open. Occasionally, head-biting occurs, but usually the shorter tortoise backs off, conceding mating rights to the victor. The behaviour is most pronounced in saddleback subspecies, which are more aggressive and have longer necks.
The prelude to mating can be very aggressive, as the male forcefully rams the female's shell with his own and nips her legs. Mounting is an awkward process and the male must stretch and tense to maintain equilibrium in a slanting position. The concave underside of the male's shell helps him to balance when straddled over the female's shell, and brings his cloacal vent (which houses the penis) closer to the female's dilated cloaca. During mating, the male vocalises with hoarse bellows and grunts, described as "rhythmic groans". This is one of the few vocalisations the tortoise makes; other noises are made during aggressive encounters, when struggling to right themselves, and hissing as they withdraw into their shells due to the forceful expulsion of air.
Egg-laying
Females journey up to several kilometres in July to November to reach nesting areas of dry, sandy coast. Nest digging is a tiring and elaborate task which may take the female several hours a day over many days to complete. It is carried out blindly using only the hind legs to dig a -deep cylindrical hole, in which the tortoise then lays up to 16 spherical, hard-shelled eggs ranging from in mass, and the size of a billiard ball. Some observations suggest that the average clutch size for domed populations (9.6 per clutch for C. porteri on Santa Cruz) is larger than that of saddlebacks (4.6 per clutch for C. duncanensis on Pinzón). The female makes a muddy plug for the nest hole out of soil mixed with urine, seals the nest by pressing down firmly with her plastron, and leaves them to be incubated by the sun. Females may lay one to four clutches per season. Temperature plays a role in the sex of the hatchlings, with lower-temperature nests producing more males and higher-temperature nests producing more females. This is related closely to incubation time, since clutches laid early incubate during the cool season and have longer incubation periods (producing more males), while eggs laid later incubate for a shorter period in the hot season (producing more females).
Early life and maturation
Young animals emerge from the nest after four to eight months and may weigh only and measure . When the young tortoises emerge from their shells, they must dig their way to the surface, which can take several weeks, though their yolk sac can sustain them up to seven months. In particularly dry conditions, the hatchlings may die underground if they are encased by hardened soil, while flooding of the nest area can drown them. Subspecies are initially indistinguishable as they all have domed carapaces. The young stay in warmer lowland areas for their first 10–15 years, encountering hazards such as falling into cracks, being crushed by falling rocks, or excessive heat stress. The Galápagos hawk was formerly the sole native predator of the tortoise hatchlings; Darwin wrote: "The young tortoises, as soon as they are hatched, fall prey in great numbers to the buzzard". The hawk is now much rarer, but introduced feral pigs, dogs, cats, and black rats have become predators of eggs and young tortoises. The adult tortoises have no natural predators apart from humans; Darwin noted: "The old ones seem generally to die from accidents, as from falling down precipices. At least several of the inhabitants told me, they had never found one dead without some such apparent cause".
Sexual maturity is reached at around 20–25 years in captivity, possibly 40 years in the wild. Life expectancy in the wild is thought to be over 100 years, making it one of the longest-lived species in the animal kingdom. Harriet, a specimen kept in Australia Zoo, was the oldest known Galápagos tortoise, having reached an estimated age of more than 170 years before her death in 2006. Chambers notes that Harriet was probably 169 years old in 2004, although media outlets claimed the greater age of 175 at death based on a less reliable timeline.
Evolutionary history
All subspecies of Galápagos tortoises evolved from common ancestors that arrived from mainland South America by overwater dispersal. Genetic studies have shown that the Chaco tortoise of Argentina and Paraguay is their closest living relative. The minimal founding population was a pregnant female or a breeding pair. Survival on the 1000-km oceanic journey is accounted for because the tortoises are buoyant, can breathe by extending their necks above the water, and are able to survive months without food or fresh water. As they are poor swimmers, the journey was probably a passive one facilitated by the Humboldt Current, which diverts westwards towards the Galápagos Islands from the mainland. The ancestors of the genus Chelonoidis are believed to have similarly dispersed from Africa to South America during the Oligocene.
The closest living relative (though not a direct ancestor) of the Galápagos giant tortoise is the Chaco tortoise (Chelonoidis chilensis), a much smaller subspecies from South America. The divergence between C. chilensis and C. niger probably occurred 11.95–25 million years ago, an evolutionary event preceding the volcanic formation of the oldest modern Galápagos Islands 5 million years ago. Mitochondrial DNA analysis indicates that the oldest existing islands (Española and San Cristóbal) were colonised first, and that these populations seeded the younger islands via dispersal in a "stepping stone" fashion via local currents. Restricted gene flow between isolated islands then resulted in the independent evolution of the populations into the divergent forms observed in the modern subspecies. The evolutionary relationships between the subspecies thus echo the volcanic history of the islands.
Darwin's development of theory of evolution
Charles Darwin visited the Galápagos for five weeks on the second voyage of HMS Beagle in 1835 and saw Galápagos tortoises on San Cristobal (Chatham) and Santiago (James) Islands. They appeared several times in his writings and journals, and played a role in the development of the theory of evolution.
Darwin wrote in his account of the voyage:
I have not as yet noticed by far the most remarkable feature in the natural history of this archipelago; it is, that the different islands to a considerable extent are inhabited by a different set of beings. My attention was first called to this fact by the Vice-Governor, Mr. Lawson, declaring that the tortoises differed from the different islands, and that he could with certainty tell from which island any one was brought ... The inhabitants, as I have said, state that they can distinguish the tortoises from the different islands; and that they differ not only in size, but in other characters. Captain Porter has described* those from Charles and from the nearest island to it, namely, Hood Island, as having their shells in front thick and turned up like a Spanish saddle, while the tortoises from James Island are rounder, blacker, and have a better taste when cooked.
The significance of the differences in tortoises between islands did not strike him as important until it was too late, as he continued,
I did not for some time pay sufficient attention to this statement, and I had already partially mingled together the collections from two of the islands. I never dreamed that islands, about fifty or sixty miles apart, and most of them in sight of each other, formed of precisely the same rocks, placed under a quite similar climate, rising to a nearly equal height, would have been differently tenanted.
Though the Beagle departed from the Galápagos with over 30 adult tortoises on deck, these were not for scientific study, but a source of fresh meat for the Pacific crossing. Their shells and bones were thrown overboard, leaving no remains with which to test any hypotheses. It has been suggested that this oversight was made because Darwin only reported seeing tortoises on San Cristóbal (C. chathamensis) and Santiago (C. darwini), both of which have an intermediate type of shell shape and are not particularly morphologically distinct from each other. Though he did visit Floreana, the C. niger subspecies found there was already nearly extinct and he was unlikely to have seen any mature animals.
However, Darwin did have four live juvenile specimens to compare from different islands. These were pet tortoises taken by himself (from San Salvador), his captain FitzRoy (two from Española) and his servant Syms Covington (from Floreana). Unfortunately, they could not help to determine whether each island had its own variety because the specimens were not mature enough to exhibit morphological differences. Although the British Museum had a few specimens, their provenance within the Galápagos was unknown. However, conversations with the naturalist Gabriel Bibron, who had seen the mature tortoises of the Paris Natural History Museum confirmed to Darwin that distinct varieties occurred.
Darwin later compared the different tortoise forms with those of mockingbirds, in the first tentative statement linking his observations from the Galapagos with the possibility of subspecies transmuting:
When I recollect the fact that [from] the form of the body, shape of scales and general size, the Spaniards can at once pronounce from which island any tortoise may have been brought; when I see these islands in sight of each other and possessed of but a scanty stock of animals, tenanted by these birds, but slightly differing in structure and filling the same place in nature; I must suspect they are only varieties ... If there is the slightest foundation for these remarks, the zoology of archipelagos will be well worth examining; for such facts would undermine the stability of subspecies.
His views on the mutability of subspecies were restated in his notebooks: "animals on separate islands ought to become different if kept long enough apart with slightly differing circumstances.—Now Galapagos Tortoises, Mocking birds, Falkland Fox, Chiloe fox,—Inglish and Irish Hare." These observations served as counterexamples to the prevailing contemporary view that subspecies were individually created.
Darwin also found these "antediluvian animals" to be a source of diversion: "I frequently got on their backs, and then giving a few raps on the hinder part of their shells, they would rise up and walk away;—but I found it very difficult to keep my balance".
Conservation
Several waves of human exploitation of the tortoises as a food source caused a decline in the total wild population from around 250,000 when first discovered in the 16th century to a low of 3,060 individuals in a 1974 census. Modern conservation efforts have subsequently brought tortoise numbers up to 19,317 (estimate for 1995–2009).
The subspecies C. n. niger became extinct by human exploitation in the 19th century. Another subspecies, C. n. abingdonii, became extinct on 24 June 2012 with the death in captivity of the last remaining specimen, a male named Lonesome George, the world's "rarest living creature". All the other surviving subspecies are listed by the IUCN as at least "vulnerable" in conservation status, if not worse.
Historical exploitation
An estimated 200,000 animals were taken before the 20th century. The relatively immobile and defenceless tortoises were collected and stored live on board ships, where they could survive for at least a year without food or water (some anecdotal reports suggest individuals surviving two years), providing valuable fresh meat, while their diluted urine and the water stored in their neck bags could be used as drinking water. The 17th-century English pirate, explorer, and naturalist William Dampier wrote, "They are so extraordinarily large and fat, and so sweet, that no pullet eats more pleasantly," while Captain James Colnett of the Royal Navy wrote of "the land tortoise which in whatever way it was dressed, was considered by all of us as the most delicious food we had ever tasted." US Navy captain David Porter declared, "after once tasting the Galapagos tortoises, every other animal food fell off greatly in our estimation ... The meat of this animal is the easiest of digestion, and a quantity of it, exceeding that of any other food, can be eaten without experiencing the slightest of inconvenience." Darwin was less enthusiastic about the meat, writing "the breast-plate roasted (as the Gauchos do "carne con cuero"), with the flesh on it, is very good; and the young tortoises make excellent soup; but otherwise the meat to my taste is indifferent."
In the 17th century, pirates started to use the Galápagos Islands as a base for resupply, restocking on food and water, and repairing vessels before attacking Spanish colonies on the South American mainland. However, the Galápagos tortoises did not struggle for survival at this point because the islands were distant from busy shipping routes and harboured few valuable natural resources. As such, they remained unclaimed by any nation, uninhabited and uncharted. In comparison, the tortoises of the islands in the Indian Ocean were already facing extinction by the late 17th century. Between the 1790s and the 1860s, whaling ships and fur sealers systematically collected tortoises in far greater numbers than the buccaneers preceding them. Some were used for food and many more were killed for high-grade "turtle oil" from the late 19th century onward for lucrative sale to continental Ecuador. A total of over 13,000 tortoises is recorded in the logs of whaling ships between 1831 and 1868, and an estimated 100,000 were taken before 1830. Since it was easiest to collect tortoises around coastal zones, females were most vulnerable to depletion during the nesting season. The collection by whalers came to a halt eventually through a combination of the scarcity of tortoises that they had created and the competition from crude oil as a cheaper energy source.
Galápagos tortoise exploitation dramatically increased with the onset of the California Gold Rush in 1849. Tortoises and sea turtles were imported into San Francisco, Sacramento and various other Gold Rush towns throughout Alta California to feed the gold mining population. Galápagos tortoise and sea turtle bones were also recovered from the Gold Rush-era archaeological site, Thompson's Cove (CA-SFR-186H), in San Francisco, California.
Population decline accelerated with the early settlement of the islands in the early 19th century, leading to unregulated hunting for meat, habitat clearance for agriculture, and the introduction of alien mammal subspecies. Feral pigs, dogs, cats, and black rats have become predators of eggs and young tortoises, whilst goats, donkeys, and cattle compete for grazing and trample nest sites. The extinction of the Floreana subspecies in the mid-19th century has been attributed to the combined pressures of hunting for the penal colony on the relatively small island, the conversion of the grazing highlands into land for farming and fruit plantations, and the introduction of feral mammals.
Scientific collection expeditions took 661 tortoises between 1888 and 1930, and more than 120 tortoises have been taken by poachers since 1990. Threats continue today with the rapid expansion of the tourist industry and increasing size of human settlements on the islands. The tortoises are down from 15 different types of subspecies when Darwin first arrived to the current 11 subspecies.
Threats
Introduced mammals
Poachers
Destruction of habitat
Characteristics that make tortoises vulnerable
Slow growth rate
Late sexual maturity
Can only be found on the Galápagos Islands
Large size and slow-moving
Collection
The tortoises of the Galápagos Islands were not only hunted for the oil that they held for fuel but also, once they were becoming more and more scarce, people began to pay to have them in their collections, as well as being put into museums.
Modern conservation
The remaining subspecies of tortoise range in IUCN classification from extinct in the wild to vulnerable. Slow growth rate, late sexual maturity, and island endemism make the tortoises particularly prone to extinction without help from conservationists. The Galápagos giant tortoise has become a flagship species for conservation efforts throughout the Galápagos.
Legal protection
The Galápagos giant tortoise is now strictly protected and is listed on Appendix I of the Convention on International Trade in Endangered subspecies of Wild Fauna and Flora (CITES). The listing requires that trade in the taxon and its products is subject to strict regulation by ratifying states, and international trade for primarily commercial purposes is prohibited. In 1936, the Ecuadorian government listed the giant tortoise as a protected subspecies. In 1959, it declared all uninhabited areas in the Galápagos to be a national park and established the Charles Darwin Foundation. In 1970, capturing or removing many subspecies from the islands (including tortoises and their eggs) was banned. To halt trade in the tortoises altogether, it became illegal to export the tortoises from Ecuador, captive or wild, continental, or insular in provenance. The banning of their exportation resulted in automatic prohibition of importation to the United States under Public Law 91-135 (1969). A 1971 Ecuadorian decree made it illegal to damage, remove, alter, or disturb any organism, rock, or other natural object in the national park.
Captive breeding
With the establishment of the Galapagos National Park and the CDF in 1959, a review of the status of the tortoise populations began. Only 11 of the 14 original populations remained and most of these were endangered if not already on the brink of extinction. The breeding and rearing program for giant tortoises began in response to the condition of the population on Pinzón, where fewer than 200 old adults were found. All of the hatchlings had been killed by introduced black rats, for perhaps more than a century. Without help, this population would eventually disappear. The only thing preserving it was the longevity of the tortoise. Its genetic resistance to the negative effects of inbreeding would be another.
Breeding and release programs began in 1965 and have successfully brought seven of the eight endangered subspecies up to less perilous population levels. Young tortoises are raised at several breeding centres across the islands to improve their survival during their vulnerable early development. Eggs are collected from threatened nesting sites, and the hatched young are given a head start by being kept in captivity for four to five years to reach a size with a much better chance of survival to adulthood, before release onto their native ranges.
The most significant population recovery was that of the Española tortoise (C. n. hoodensis), which was saved from near-certain extinction. The population had been depleted to three males and 12 females that had been so widely dispersed that no mating in the wild had occurred. Fruitless attempts to breed one of the tortoises, Lonesome George for example, is speculated to be attributed to a lack of postnatal cues, and confusion over which would be the most appropriate genetic subspecies would be the most appropriate to mate him with on the islands. The 15 remaining tortoises were brought to the Charles Darwin Research Station in 1971 for a captive breeding program and, in the following 33 years, they gave rise to over 1,200 progeny which were released onto their home island and have since begun to reproduce naturally. One of the tortoises, Diego, is one of the main drivers of a remarkable recovery of the hoodensis subspecies, having fathered between 350 and 800 progeny.
Island restoration
The Galápagos National Park Service systematically culls feral predators and competitors. Goat eradication on islands, including Pinta, was achieved by the technique of using "Judas goats" with radio location collars to find the herds. Marksmen then shot all the goats except the Judas, and then returned weeks later to find the "Judas" and shoot the herd to which it had relocated. Goats were removed from Pinta Island after a 30-year eradication campaign, the largest removal of an insular goat population using ground-based methods. Over 41,000 goats were removed during the initial hunting effort (1971–82). This process was repeated until only the "Judas" goat remained, which was then killed. Other measures have included dog eradication from San Cristóbal, and fencing off nests to protect them from feral pigs.
Efforts are now underway to repopulate islands formerly inhabited by tortoises to restore their ecosystems (island restoration) to their condition before humans arrived. The tortoises are a keystone species, acting as ecosystem engineers which help in plant seed dispersal and trampling down brush and thinning the understory of vegetation (allowing light to penetrate and germination to occur). Birds such as flycatchers perch on and fly around tortoises to hunt the insects they displace from the brush. In May 2010, 39 sterilised tortoises of hybrid origin were introduced to Pinta Island, the first tortoises there since the evacuation of Lonesome George 38 years before. Sterile tortoises were released so the problem of interbreeding between subspecies would be avoided if any fertile tortoises were to be released in the future. It is hoped that with the recent identification of a hybrid C. n. abingdonii tortoise, the approximate genetic constitution of the original inhabitants of Pinta may eventually be restored with the identification and relocation of appropriate specimens to this island. This approach may be used to "retortoise" Floreana in the future, since captive individuals have been found to be descended from the extinct original stock.
Applied science
The Galapagos Tortoise Movement Ecology Programme is a collaborative project coordinated by Dr Stephen Blake of the Max Planck Institute for Ornithology. Its goal is to assist the Galapagos National Park to effectively conserve giant tortoises by conducting cutting-edge applied science, and developing an inspirational tortoise-based outreach and education programme. Since 2009, the project team have been analysing the movements of giant tortoises by tracking them via satellite tags. As of November 2014, the team have tagged 83 tortoises from four subspecies on three islands. They have established that giant tortoises conduct migrations up and down volcanoes, primarily in response to seasonal changes in the availability and quality of vegetation. In 2015 they will start to track the movements of hatchling and juvenile tortoises, supported by the UK's Galapagos Conservation Trust.
| Biology and health sciences | Reptiles | null |
7936253 | https://en.wikipedia.org/wiki/Red%20king%20crab | Red king crab | The red king crab (Paralithodes camtschaticus), also called Kamchatka crab or Alaskan king crab, is a species of king crab native to cold waters in the North Pacific Ocean and adjacent seas, but also introduced to the Barents Sea. It grows to a leg span of , and is heavily targeted by fisheries.
Description
The red king crab is the largest species of king crab. Red king crabs can reach a carapace width up to , a leg span of , and a weight of . Males grow larger than females. Today, red king crabs infrequently surpass in carapace width and the average male landed in the Bering Sea weighs . It was named after the color it turns when it is cooked rather than the color of a living animal, which tends to be more burgundy.
Distribution
The red king crab is native to cold waters in the North Pacific Ocean and adjacent seas, ranging from the Bering Sea south to the Gulf of Alaska, off the Kamchatka Peninsula, and in the Sea of Okhotsk and Sea of Japan. It was introduced artificially by the Soviet Union into the Murmansk Fjord, Barents Sea, during the 1960s to provide a new, and valuable, catch in Europe.
Red king crabs have been seen in water temperatures that range from , with typical being . Immatures prefer temperatures below . The depth at which it can live has much to do with what stage of its lifecycle it is in; newly hatched crab (zoea larvae) stay in the shallower waters where food and protection are plentiful. Usually, after the age of two, the crabs move down to depths of and take part in what is known as podding; hundreds of crabs come together in tight, highly concentrated groups. Adult crabs are found usually more than down on the sand and muddy areas in the substrate. They migrate in the winter or early spring to shallower depths for mating, but most of their lives are spent in the deep waters where they feed.
Ecology
P. camtschaticus faces many predators in its native range including Pacific cod, walleye pollock, rock sole, flathead sole, rex sole, Dover sole (Microstomus pacificus), arrowtooth flounder, Elasmobranchs, halibut, sculpin, Greenland turbot, Pacific salmon, Pacific herring, otters (Enhydra lutris) and seals.
Fisheries
The red king crab is the most coveted of the commercially sold king crab species, and it is the most expensive per unit weight. It is most commonly caught in the Bering Sea and Norton Sound, Alaska, and is particularly difficult to catch, but is nonetheless one of the most preferred crabs for consumption.
Red king crabs are experiencing a steady decline in numbers in their native far east coastal waters for unclear reasons, though several theories for the precipitous drop in the crab population have been proposed, including overfishing, warmer waters, and increased fish predation. Fishing controls set by the United States in the 1980s and 2000s have failed to stem the decline.
In Europe
In the 1960s, the Soviet Union transported red king crabs from the North Pacific Ocean to the Murmansk Fjord. They did not survive transport overland, so a batch was flown in, which survived, was released, and bred and spread in the wild.
It was first found in Norway in 1977. In the Barents Sea, it is an invasive species and its population is increasing tremendously. This is causing great concern to local environmentalists and local fishermen, as the crab eats everything it comes across and is spreading very rapidly, eating most seabed life and "creating a desert". Since its introduction in the 1960s, it has spread west along the Norwegian coast and also has reached Svalbard. The species keeps on advancing southwards along the coast of Norway and some scientists think they are advancing around a year. In Norway they are sometimes called "Stalin's crabs", since they were introduced by the Soviet Union.
By the mid 1990s, the king crabs reached North Cape. The Norwegian Institute of Marine Research found in 2010–2013 that they have reached Sørøya and are breeding there. A few have been caught as far south as Tromsø. There is fear of the result if they reach the cod breeding grounds off Lofoten.
A report on 8 June 2009 said that a red king crab had been caught off Skogsvåg at Sotra south of Bergen in south Norway. An important natural predator of the red king crab, the giant Pacific octopus (Enteroctopus dofleini, formerly called Octopus apollyon), does not occur in European waters. A fisherman in Honningsvåg (a town near the North Cape) complained that king crabs' claws were ruining fishing nets and deep lines. Despite these concerns, the species is protected by diplomatic accords between Norway and Russia as part of a fisheries agreement between Norway and Russia about the Barents Sea, and a bilateral fishing commission decides how to manage the stocks and imposes fishing quotas. West of the North Cape on Norway's northern tip, Norway manages its crab population itself. As of May 2006, only 259 Norwegian fishermen were allowed to catch it east of the North Cape.
In the Norwegian Sea, some evidence indicates that the red king crabs eat the egg masses of the capelin, which is an important prey for the cod. The report (as at 24 May 2006) said that in Norwegian Sea, in the Barents Sea (east of the North Cape), catching red king crab is allowed with license only due to a fisheries agreement between Norway and Russia, but elsewhere in Norwegian seas, the catching of king crab is much freer, but nonetheless, if someone catches one, it is illegal to throw it back in the sea.
In January 2022 it was reported that fishermen in the United Kingdom had caught red king crabs, but they were later identified as the native Lithodes maja.
Waste recycling
On average, crab processing waste can account for 69% of the catch mass. The mass fraction of carapace from these wastes is approximately 60%; the rest comprises the entrails (including the digestive organ, the hepatopancreas). In red king crab, the hepatopancreas makes up about 90% of the intestines of the carapace and 5–10% of the total weight of the animal. The hepatopancreas of the digestive system of commercial crabs is a valuable source of a complex of enzymes with various activities: collagenase, protease, hyaluronidase, lipase, nuclease, etc. The complex of proteolytic enzymes of the red king crab hepatopancreas is of interest in various industries.
Physiology
Mature female red king crabs must stay in warmer water (near ) to ensure the eggs will be ready for hatching, while the male red king crabs stay in relatively cold water (near ) to conserve energy. In spring (May), female red king crabs move to shallow coastal areas to molt and spawn, and males join the females in the shallow water before molting. In the summer (mid-June through mid-November), these crabs spend their time in fairly deep water, below the established summer thermocline. When the thermocline breaks down, the red king crabs migrate back to intermediate depths, where they stay until the female red king crabs release the eggs fertilized in the previous spawning.
The red king crab has a wide range of tolerance to temperature, but it affects their growth. The organism's growth and molting is slow when outside temperature falls below ; around , they molt rather quickly.
Overall, red king crabs have a high adaptation capacity in changes of salinity level because the crabs retain their vital functions and their feeding activities. A difference is seen, though, in the salinity tolerance between juvenile and adult red king crabs. Juveniles are slightly more tolerant to low salinity because their volume regulation is significantly better. Juveniles are consistently hyposmotic to the seawater because they have lower sodium concentration in their hemolymph. As the juveniles are smaller, their exoskeleton is more rigid. The adult red king crabs are hyperosmotic in high salinity and becomes hyposmotic in lower salinity. The hyperosmoticity is due to the higher sodium and potassium concentrations in the hemolymph compared to the surrounding water they live in.
A slight fluctuation on the pH level of the water (i.e. making the water more acidic) would have great effect on the red king crab. They grow slower in acidified water (pH 7.8 instead of 8.0) and eventually die after longer exposure times because of the imbalance of the organisms' acid-base equilibrium.
Respiration
The red king crab has five sets of gills used for respiration, which are in the bronchial chamber within the carapace. The carapace is a covering of sheets of exoskeleton that overhang the thorax vertically to fit over the base of the thoracic legs. The carapace encloses two branchial chambers that enclose the gills. The gill surfaces are covered in chitinous cuticle, which is permeable to gases, allowing gas exchange. Internal gills, like other specialized gills, need metabolic energy to pull water over the respiratory surface. To induce a current into the branchial chamber the crab uses back and forth movements of an appendage called the scaphognathite. The water is drawn in from behind the walking legs then expelled from the branchial chambers through the tubes called prebronchial apertures, which are located beside the mouth. To filter the water before entering the branchial chamber, crabs have branchiostegal hairs that can collect debris. Due to the environment to which it is exposed, the posterior gills of the crab can also be cleared of parasites and sediment by increasing the movement of its fifth set of primitive legs.
Each gill has a main axis with many lateral filaments or lamellae that are vascularized. The afferent channel transports blood from the gill axis into each filament through a fine afferent canal to the gill top. Blood returns by a minute efferent canal to the gill tip to the efferent channel and passes to the pericardial chamber, which contains the heart. Gases are exchanged in the numerous filaments, and oxygen absorption is especially facilitated by hemocyanin. Red king crabs exhibit unidirectional ventilation. This can be described as the flow of water in a U-shaped course; water passes posteriorly from the incurrent opening, an opening in the carapace near the base of the chelipeds, dorsally over the gills, and anteriorly to exit beside the head.
Circulation
Due to their respiratory system's limited ability to deliver by diffusion, respiratory gases must be transported around the body.Paralithodes camtschaticus has an open circulatory system with a dorsal, ostiate heart. An open circulatory system has circulating fluid that passes somewhat freely among the tissues before being collected and recirculated. The heart is in a pericardial chamber, and blood passes through this chamber into the lumen of the heart through two pairs of ostia. Seven arteries conduct blood from the heart to various regions of the body. Each artery branches extensively, and smaller arteries ultimately end in the hemocoel. Venous blood drains into the sternal sinus, where it is conveyed by channels to the gills for aeration and returned again to the pericardial sinus.
They have a neurogenic heart, which has rhythmic depolarization that is responsible for initiating heartbeats. Heartbeats originate in nervous tissue; innervated muscle cells cause the heart to contract when stimulated by nerve impulses. The cardiac ganglion, which consists of nine neurons, attaches to the dorsal wall of the heart. The anterior neurons innervate the heart, whereas the other posterior neurons make synaptic contact with those anterior neurons. The posterior neuron acts as the pacemaker but also functions as the cellular oscillator and the central pattern generator. This posterior neuron produces a train of impulses, which excites the other posterior neurons. The heart contracts when the posterior neurons activate the five anterior neurons, which send impulses to the muscle cells. This is how the Frank–Starling mechanism works within crustaceans. The Frank-Starling mechanism refers to the vitally important intrinsic control of the heart; mainly, the stretching of the cardiac muscle tends to increase the force of its contraction by an effect at the cellular level. This mechanism is important as it allows the organism to match its output of blood with its input of blood. Because of the stretching between beats, the Frank-Starling mechanism allows the heart to then naturally contract more forcefully, allowing greater flow of blood, which results in the matched heart output to the increased blood received. The Frank-Starling mechanism is a little different in crustaceans, as it involves the cardiac ganglion as described previously. The stretching of the heart induces the ganglion to fire more regularly and powerfully.
Red king crab blood contains leukocytes and the second-most common respiratory pigment called hemocyanin. Arthropod hemocyanin is a distinct variation specific to arthropods and is a metalloprotein that uses copper atoms that are bound to its structure. Two copper atoms are needed to bind one O2 molecule. Because it is a large protein molecule, it is found in the blood plasma, but not in body tissues or muscles. Hemocyanins are named appropriately because when oxygenated, their color changes from colorless to blue.
| Biology and health sciences | Crabs and hermit crabs | Animals |
7937417 | https://en.wikipedia.org/wiki/Pantodonta | Pantodonta | Pantodonta is an extinct suborder (or, according to some, an order) of eutherian mammals. These herbivorous mammals were one of the first groups of large mammals to evolve (around 66 million years ago) after the end of the Cretaceous. The last pantodonts died out at the end of the Eocene (around 34 million years ago).
Pantodonta include some of the largest mammals of their time, but were a diversified group, with some primitive members weighing less than and the largest more than .
The earliest and most primitive pantodonts, Bemalambda (with a skull probably the size of a dog) and Hypsilolambda, appear in the early Paleocene Shanghuan Formation in China. All more derived families are collectively classified as Eupantodonta. The pantodonts appear in North America in the middle Paleocene, where Coryphodon survived into the middle Eocene. Pantodont teeth have been found in South America (Alcidedorbignya) and Antarctica, and footprints in a coal mine on Svalbard.
Description
The pantodonts varied considerably in size: the small Archaeolambda, of which there is a complete skeleton from the Late Palaeocene of China, was probably arboreal, while the North American, ground sloth-like Barylambda was massive, slow-moving ("graviportal") and probably browsed on high vegetation.
Dentition
The pantodonts have a primitive dental formula () with little or no diastemata. Their most important synapomorphy are the zalambdodont (V-shaped ectoloph opening towards lip) P3–4 and (except in the most primitive families) dilambdodont (W-shaped ectoloph) upper molars. Most pantodonts lacked a hypocone (fourth cusp) and had small conules (additional small cusps). The incisors are small but the canines large, occasionally sabertooth-like. On P3-M3 there is normally an ectoflexus (indentation on the outer side). Asian families can typically be distinguished from the American because their paracone and metacone (bottom of W on side of tongue) tend to be closer together.
The cheek teeth in the lower jaw are also dilambdodont, with broad, high metalophids (posterior crest) and tall metaconid (posterior-interior cusp) with much lower paracristids and small paraconids.
Postcranial skeleton
Pantodonts have plesiomorphic (unaltered) and robust postcranial skeletons. Their five-toed feet are often hoofed with the tarsals similar to those of ungulates, which feature had led to previously suggested ties to arctocyonid "condylarths", but this similarity is now considered primitive.
Classification
The pantodonts were previously grouped with the ungulates as amblypods, paenungulates, or arctocyonids, but since they have been allied with the tillodonts and considered to be derived from the cimolestids. The interrelationship within Pantodonta is controversial, but, following , it contains about two dozen genera in ten families. Most of the families are known from the Paleocene of either Asia or North America. The pantolambdodontids and coryphodontids survived into the Eocene and the latter are known from across the northern hemisphere. Some dental features can possibly link the most primitive pantodonts to the palaeoryctids, a group of small and insectivorous mammals that evolved during the Cretaceous. Recently a close relationship with Periptychidae has been suggested. This would make pantodonts crown-group ungulate placentals and not related to cimolestids at all.
Genera from North America tended to be large and robust, starting with Pantolambda and Caenolambda in the Middle Paleocene epoch, and later in the epoch started to get larger, with Barylambda as the largest Paleocene form of pantodont. However, Asian forms, such as Archaeolambda, tended to be thinner and less robust, around the size of a medium-sized dog. Only later in the Eocene, with Hypercoryphodon, did Asian pantodonts get large and robust.
Timeline of genera
| Biology and health sciences | Mammals: General | Animals |
429425 | https://en.wikipedia.org/wiki/Probability%20amplitude | Probability amplitude | In quantum mechanics, a probability amplitude is a complex number used for describing the behaviour of systems. The square of the modulus of this quantity represents a probability density.
Probability amplitudes provide a relationship between the quantum state vector of a system and the results of observations of that system, a link was first proposed by Max Born, in 1926. Interpretation of values of a wave function as the probability amplitude is a pillar of the Copenhagen interpretation of quantum mechanics. In fact, the properties of the space of wave functions were being used to make physical predictions (such as emissions from atoms being at certain discrete energies) before any physical interpretation of a particular function was offered. Born was awarded half of the 1954 Nobel Prize in Physics for this understanding, and the probability thus calculated is sometimes called the "Born probability". These probabilistic concepts, namely the probability density and quantum measurements, were vigorously contested at the time by the original physicists working on the theory, such as Schrödinger and Einstein. It is the source of the mysterious consequences and philosophical difficulties in the interpretations of quantum mechanics—topics that continue to be debated even today.
Physical overview
Neglecting some technical complexities, the problem of quantum measurement is the behaviour of a quantum state, for which the value of the observable to be measured is uncertain. Such a state is thought to be a coherent superposition of the observable's eigenstates, states on which the value of the observable is uniquely defined, for different possible values of the observable.
When a measurement of is made, the system (under the Copenhagen interpretation) jumps to one of the eigenstates, returning the eigenvalue belonging to that eigenstate. The system may always be described by a linear combination or superposition of these eigenstates with unequal "weights". Intuitively it is clear that eigenstates with heavier "weights" are more "likely" to be produced. Indeed, which of the above eigenstates the system jumps to is given by a probabilistic law: the probability of the system jumping to the state is proportional to the absolute value of the corresponding numerical weight squared. These numerical weights are called probability amplitudes, and this relationship used to calculate probabilities from given pure quantum states (such as wave functions) is called the Born rule.
Clearly, the sum of the probabilities, which equals the sum of the absolute squares of the probability amplitudes, must equal 1. This is the normalization requirement.
If the system is known to be in some eigenstate of (e.g. after an observation of the corresponding eigenvalue of ) the probability of observing that eigenvalue becomes equal to 1 (certain) for all subsequent measurements of (so long as no other important forces act between the measurements). In other words, the probability amplitudes are zero for all the other eigenstates, and remain zero for the future measurements. If the set of eigenstates to which the system can jump upon measurement of is the same as the set of eigenstates for measurement of , then subsequent measurements of either or always produce the same values with probability of 1, no matter the order in which they are applied. The probability amplitudes are unaffected by either measurement, and the observables are said to commute.
By contrast, if the eigenstates of and are different, then measurement of produces a jump to a state that is not an eigenstate of . Therefore, if the system is known to be in some eigenstate of (all probability amplitudes zero except for one eigenstate), then when is observed the probability amplitudes are changed. A second, subsequent observation of no longer certainly produces the eigenvalue corresponding to the starting state. In other words, the probability amplitudes for the second measurement of depend on whether it comes before or after a measurement of , and the two observables do not commute.
Mathematical formulation
In a formal setup, the state of an isolated physical system in quantum mechanics is represented, at a fixed time , by a state vector belonging to a separable complex Hilbert space. Using bra–ket notation the relation between state vector and "position basis" of the Hilbert space can be written as
.
Its relation with an observable can be elucidated by generalizing the quantum state to a measurable function and its domain of definition to a given -finite measure space . This allows for a refinement of Lebesgue's decomposition theorem, decomposing μ into three mutually singular parts
where μac is absolutely continuous with respect to the Lebesgue measure, μsc is singular with respect to the Lebesgue measure and atomless, and μpp is a pure point measure.
Continuous amplitudes
A usual presentation of the probability amplitude is that of a wave function belonging to the space of (equivalence classes of) square integrable functions, i.e., belongs to if and only if
.
If the norm is equal to and such that
,
then is the probability density function for a measurement of the particle's position at a given time, defined as the Radon–Nikodym derivative with respect to the Lebesgue measure (e.g. on the set of all real numbers). As probability is a dimensionless quantity, must have the inverse dimension of the variable of integration . For example, the above amplitude has dimension [L−1/2], where L represents length.
Whereas a Hilbert space is separable if and only if it admits a countable orthonormal basis, the range of a continuous random variable is an uncountable set (i.e. the probability that the system is "at position " will always be zero). As such, eigenstates of an observable need not necessarily be measurable functions belonging to (see normalization condition below). A typical example is the position operator defined as
whose eigenfunctions are Dirac delta functions
which clearly do not belong to . By replacing the state space by a suitable rigged Hilbert space, however, the rigorous notion of eigenstates from spectral theorem as well as spectral decomposition is preserved.
Discrete amplitudes
Let be atomic (i.e. the set in is an atom); specifying the measure of any discrete variable equal to . The amplitudes are composed of state vector indexed by ; its components are denoted by for uniformity with the previous case. If the -norm of is equal to 1, then is a probability mass function.
A convenient configuration space is such that each point produces some unique value of the observable . For discrete it means that all elements of the standard basis are eigenvectors of . Then is the probability amplitude for the eigenstate . If it corresponds to a non-degenerate eigenvalue of , then gives the probability of the corresponding value of for the initial state .
if and only if is the same quantum state as . if and only if and are orthogonal. Otherwise the modulus of is between 0 and 1.
A discrete probability amplitude may be considered as a fundamental frequency in the probability frequency domain (spherical harmonics) for the purposes of simplifying M-theory transformation calculations. Discrete dynamical variables are used in such problems as a particle in an idealized reflective box and quantum harmonic oscillator.
Examples
An example of the discrete case is a quantum system that can be in two possible states, e.g. the polarization of a photon. When the polarization is measured, it could be the horizontal state or the vertical state . Until its polarization is measured the photon can be in a superposition of both these states, so its state could be written as
,
with and the probability amplitudes for the states and respectively. When the photon's polarization is measured, the resulting state is either horizontal or vertical. But in a random experiment, the probability of being horizontally polarized is , and the probability of being vertically polarized is .
Hence, a photon in a state would have a probability of to come out horizontally polarized, and a probability of to come out vertically polarized when an ensemble of measurements are made. The order of such results, is, however, completely random.
Another example is quantum spin. If a spin-measuring apparatus is pointing along the z-axis and is therefore able to measure the z-component of the spin (), the following must be true for the measurement of spin "up" and "down":
If one assumes that system is prepared, so that +1 is registered in and then the apparatus is rotated to measure , the following holds:
The probability amplitude of measuring spin up is given by , since the system had the initial state . The probability of measuring is given by
Which agrees with experiment.
Normalization
In the example above, the measurement must give either or , so the total probability of measuring or must be 1. This leads to a constraint that ; more generally the sum of the squared moduli of the probability amplitudes of all the possible states is equal to one. If to understand "all the possible states" as an orthonormal basis, that makes sense in the discrete case, then this condition is the same as the norm-1 condition explained above.
One can always divide any non-zero element of a Hilbert space by its norm and obtain a normalized state vector. Not every wave function belongs to the Hilbert space , though. Wave functions that fulfill this constraint are called normalizable.
The Schrödinger equation, describing states of quantum particles, has solutions that describe a system and determine precisely how the state changes with time. Suppose a wave function gives a description of the particle (position at a given time ). A wave function is square integrable if
After normalization the wave function still represents the same state and is therefore equal by definition to
Under the standard Copenhagen interpretation, the normalized wavefunction gives probability amplitudes for the position of the particle. Hence, is a probability density function and the probability that the particle is in the volume at fixed time is given by
The probability density function does not vary with time as the evolution of the wave function is dictated by the Schrödinger equation and is therefore entirely deterministic. This is key to understanding the importance of this interpretation: for a given particle constant mass, initial and potential, the Schrödinger equation fully determines subsequent wavefunctions. The above then gives probabilities of locations of the particle at all subsequent times.
In the context of the double-slit experiment
Probability amplitudes have special significance because they act in quantum mechanics as the equivalent of conventional probabilities, with many analogous laws, as described above. For example, in the classic double-slit experiment, electrons are fired randomly at two slits, and the probability distribution of detecting electrons at all parts on a large screen placed behind the slits, is questioned. An intuitive answer is that , where is the probability of that event. This is obvious if one assumes that an electron passes through either slit. When no measurement apparatus that determines through which slit the electrons travel is installed, the observed probability distribution on the screen reflects the interference pattern that is common with light waves. If one assumes the above law to be true, then this pattern cannot be explained. The particles cannot be said to go through either slit and the simple explanation does not work. The correct explanation is, however, by the association of probability amplitudes to each event. The complex amplitudes which represent the electron passing each slit ( and ) follow the law of precisely the form expected: . This is the principle of quantum superposition. The probability, which is the modulus squared of the probability amplitude, then, follows the interference pattern under the requirement that amplitudes are complex:
Here, and
are the arguments of and respectively. A purely real formulation has too few dimensions to describe the system's state when superposition is taken into account. That is, without the arguments of the amplitudes, we cannot describe the phase-dependent interference. The crucial term is called the "interference term", and this would be missing if we had added the probabilities.
However, one may choose to devise an experiment in which the experimenter observes which slit each electron goes through. Then, due to wavefunction collapse, the interference pattern is not observed on the screen.
One may go further in devising an experiment in which the experimenter gets rid of this "which-path information" by a "quantum eraser". Then, according to the Copenhagen interpretation, the case A applies again and the interference pattern is restored.
Conservation of probabilities and the continuity equation
Intuitively, since a normalised wave function stays normalised while evolving according to the wave equation, there will be a relationship between the change in the probability density of the particle's position and the change in the amplitude at these positions.
Define the probability current (or flux) as
measured in units of (probability)/(area × time).
Then the current satisfies the equation
The probability density is , this equation is exactly the continuity equation, appearing in many situations in physics where we need to describe the local conservation of quantities. The best example is in classical electrodynamics, where corresponds to current density corresponding to electric charge, and the density is the charge-density. The corresponding continuity equation describes the local conservation of charges.
Composite systems
For two quantum systems with spaces and and given states and respectively, their combined state can be expressed as a function on , that gives the
product of respective probability measures. In other words, amplitudes of a non-entangled composite state are products of original amplitudes, and respective observables on the systems 1 and 2 behave on these states as independent random variables. This strengthens the probabilistic interpretation explicated above .
Amplitudes in operators
The concept of amplitudes is also used in the context of scattering theory, notably in the form of S-matrices. Whereas moduli of vector components squared, for a given vector, give a fixed probability distribution, moduli of matrix elements squared are interpreted as transition probabilities just as in a random process. Like a finite-dimensional unit vector specifies a finite probability distribution, a finite-dimensional unitary matrix specifies transition probabilities between a finite number of states.
The "transitional" interpretation may be applied to s on non-discrete spaces as well.
| Physical sciences | Quantum mechanics | Physics |
429528 | https://en.wikipedia.org/wiki/Vole | Vole | Voles are small rodents that are relatives of lemmings and hamsters, but with a stouter body; a longer, hairy tail; a slightly rounder head; smaller eyes and ears; and differently formed molars (high-crowned with angular cusps instead of low-crowned with rounded cusps). They are sometimes known as meadow mice or field mice.
Vole species form the subfamily Arvicolinae with the lemmings and the muskrats. There are approximately 155 different vole species.
Description
Voles are small rodents that grow to , depending on the species. Females can have five to ten litters per year, though with an average lifespan of three months and requiring one month to adulthood, two litters is the norm. Gestation lasts for three weeks and the young voles reach sexual maturity in a month. As a result of this biological exponential growth, vole populations can grow very large within a short time. One mating pair can produce 100 offspring in a year.
Voles outwardly resemble several other small animals. Moles, gophers, mice, rats and even shrews have similar characteristics and behavioral tendencies.
Voles thrive on small plants yet, like shrews, they will eat dead animals and, like mice and rats, they can live on almost any nut or fruit. In addition, voles target plants more than most other small animals, making their presence evident. Voles readily girdle small trees and ground cover much like a porcupine. This girdling can easily kill young plants and is not healthy for trees and other shrubs.
Voles often eat succulent root systems and burrow under plants and eat away until the plant is dead. Bulbs are another favorite target for voles; their excellent burrowing and tunnelling skills give them access to sensitive areas without clear or early warning. The presence of large numbers of voles is often identifiable only after they have destroyed a number of plants. However, like other burrowing rodents, they also play beneficial roles, including dispersing nutrients throughout the upper soil layers.
Predators
Many predators eat voles, including martens, owls, hawks, falcons, coyotes, bobcats, foxes, raccoons, squirrels, snakes, weasels, domestic cats and lynxes. Vole bones are often found in the pellets of the short-eared owl, the northern spotted owl, the saw-whet owl, the barn owl, the great gray owl, and the northern pygmy owl. In the summer of 2024, biologists and other scientists at UC Davis first observed California ground squirrels actively hunting voles.
Lifespan
The average lifespan for smaller species of vole is three to six months, and they rarely live longer than 12 months. Larger species, such as the European water vole, live longer and usually die during their second, or rarely their third, winter. As many as 88% of voles are estimated to die within the first month of life.
Genetics and sexual behavior
The prairie vole is a notable animal model for its monogamous social fidelity, since the male is usually socially faithful to the female, and shares in the raising of pups. The woodland vole is also usually monogamous. Another species from the same genus, the meadow vole, has promiscuously mating males, and scientists have changed adult male meadow voles' behavior to resemble that of prairie voles in experiments in which a viral vector was used to increase a single gene's expression within a particular brain region.
The behavior is influenced by the number of repetitions of a particular string of microsatellite DNA. Male prairie voles with the longest DNA strings spend more time with their mates and pups than male prairie voles with shorter strings. However, other scientists have disputed the gene's relationship to monogamy, and cast doubt on whether the human version plays an analogous role. Physiologically, pair-bonding behavior has been shown to be connected to vasopressin, dopamine, and oxytocin levels, with the genetic influence apparently arising via the number of receptors for these substances in the brain; the pair-bonding behavior has also been shown in experiments to be strongly modifiable by administering some of these substances directly.
Voles have a number of unusual chromosomal traits. Species have been found with 17 to 64 chromosomes. In some species, males and females have different chromosome numbers, a trait unusual in mammals, though it is seen in other organisms. Additionally, genetic material typically found on the Y chromosome has been found in both males and females in at least one species. In another species, the X chromosome contains 20% of the genome. All of these variations result in very little physical aberration; most vole species are virtually indistinguishable. In one species, the creeping vole Microtus oregoni, it was discovered the Y chromosome has been lost entirely; the male-determining chromosome is actually a second X that is largely identical to the female X, and both the maternally inherited and male-specific sex chromosomes carry vestiges of the ancestral Y. This is quite unusual in mammals, as the XY system is fairly stable across a number of mammal species.
Mating system
Voles may be either monogamous or polygamous, which leads to differing patterns of mate choice and parental care. Environmental conditions play a large part in dictating which system is active in a given population. Voles live in colonies due to the young remaining in the family group for relatively long periods. In the genus Microtus, monogamy is preferred when resources are spatially homogeneous and population densities are low; where the opposites of both conditions are realized, polygamous tendencies arise. Vole mating systems are also sensitive to the operational sex ratio and tend toward monogamy when males and females are present in equal numbers. Where one sex is more numerous than the other, polygamy is more likely. However the most marked effect on mating system is population density and these effects can take place both inter and intra-specifically.
Male voles are territorial and tend to include territories of several female voles when possible. Under these conditions polygyny exists and males offer little parental care. Males mark and aggressively defend their territories since females prefer males with the most recent marking in a given area.
Voles prefer familiar mates through olfactory sensory exploitation. Monogamous voles prefer males who have yet to mate, while non-monogamous voles do not. Mate preference in voles develops through cohabitation in as little as 24 hours. This drives young male voles to show non-limiting preference toward female siblings. This is not inclusive to females' preference for males which may help to explain the absence of interbreeding indicators.
Although females show little territoriality, under pair bonding conditions they tend to show aggression toward other female voles. This behavior is flexible as some Microtus females share dens during the winter months, perhaps to conserve heat and energy. Populations which are monogamous show relatively minor size differences between genders compared with those using polygamous systems.
The grey-sided vole (Myodes rufocanus) exhibits male-biased dispersal as a means of avoiding incestuous matings. Among those matings that involve inbreeding, the number of weaned juveniles in litters is significantly fewer than that from noninbred litters, due to inbreeding depression.
Brandt's vole (Lasiopodomys brandtii) lives in groups that mainly consist of close relatives. However, they show no sign of inbreeding. The mating system of these voles involves a type of polygyny for males and extra-group polyandry for females. This system increases the frequency of mating among distantly related individuals, and is achieved mainly by dispersal during the mating season. Such a strategy is likely an adaptation to avoid the inbreeding depression that would be caused by expression of deleterious recessive alleles if close relatives mated.
Empathy and consolation
A 2016 study into the behavior of voles, Microtus ochrogaster specifically, found that voles comfort each other when mistreated, spending more time grooming a mistreated vole. Voles that were not mistreated had levels of stress hormones that were similar to the voles that had been mistreated, suggesting that the voles were capable of empathizing with each other. This was further proven by blocking the vole's receptors for oxytocin, a hormone involved in empathy. When the oxytocin receptors were blocked this behavior stopped.
This type of empathetic behavior has previously been thought to occur only in animals with advanced cognition such as humans, apes, and elephants.
Vole clock
The vole clock is a method of dating archaeological strata using vole teeth.
Classification
Order Rodentia
Superfamily Muroidea
Family Cricetidae
Subfamily Arvicolinae (in part)
Tribe Arvicolini
Genus Arvicola – water voles
Genus Blanfordimys – Afghan vole and Bucharian vole
Genus Chionomys – snow voles
Genus Lasiopodomys
Genus Lemmiscus – sagebrush vole
Genus Microtus – voles
Genus Neodon – mountain voles
Genus Phaiomys
Genus Proedromys – Duke of Bedford's vole
Genus Volemys
Tribe Ellobiusini – mole voles
Genus Ellobius – mole voles
Tribe Clethrionomyini
Genus Alticola – voles from Central Asia
Genus Caryomys
Genus Eothenomys – voles from East Asia
Genus Hyperacrius – voles from India (True's vole), Afghanistan and Pakistan (Murree's vole)
Genus Craseomys – some red-backed voles
Genus Clethrionomys - other red-backed voles
Tribe Phenacomyini
Genus Arborimus – tree voles
Genus Phenacomys – heather voles
Tribe Pliomyini
Genus Dinaromys'' – voles from the Dinaric Alps
| Biology and health sciences | Rodents | Animals |
429542 | https://en.wikipedia.org/wiki/Preterm%20birth | Preterm birth | Preterm birth, also known as premature birth, is the birth of a baby at fewer than 37 weeks gestational age, as opposed to full-term delivery at approximately 40 weeks. Extreme preterm is less than 28 weeks, very early preterm birth is between 28 and 32 weeks, early preterm birth occurs between 32 and 34 weeks, late preterm birth is between 34 and 36 weeks' gestation. These babies are also known as premature babies or colloquially preemies (American English) or premmies (Australian English). Symptoms of preterm labor include uterine contractions which occur more often than every ten minutes and/or the leaking of fluid from the vagina before 37 weeks. Premature infants are at greater risk for cerebral palsy, delays in development, hearing problems and problems with their vision. The earlier a baby is born, the greater these risks will be.
The cause of spontaneous preterm birth is often not known. Risk factors include diabetes, high blood pressure, multiple gestation (being pregnant with more than one baby), being either obese or underweight, vaginal infections, air pollution exposure, tobacco smoking, and psychological stress. For a healthy pregnancy, medical induction of labor or cesarean section are not recommended before 39 weeks unless required for other medical reasons. There may be certain medical reasons for early delivery such as preeclampsia.
Preterm birth may be prevented in those at risk if the hormone progesterone is taken during pregnancy. Evidence does not support the usefulness of bed rest. It is estimated that at least 75% of preterm infants would survive with appropriate treatment, and the survival rate is highest among the infants born the latest in gestation. In women who might deliver between 24 and 37 weeks, corticosteroid treatment may improve outcomes. A number of medications, including nifedipine, may delay delivery so that a mother can be moved to where more medical care is available and the corticosteroids have a greater chance to work. Once the baby is born, care includes keeping the baby warm through skin-to-skin contact or incubation, supporting breastfeeding and/or formula feeding, treating infections, and supporting breathing. Preterm babies sometimes require intubation.
Preterm birth is the most common cause of death among infants worldwide. About 15 million babies are preterm each year (5% to 18% of all deliveries). Late preterm birth accounts for 75% of all preterm births. This rate is inconsistent across countries. In the United Kingdom 7.9% of babies are born pre-term and in the United States 12.3% of all births are before 37 weeks gestation. Approximately 0.5% of births are extremely early periviable births (20–25 weeks of gestation), and these account for most of the deaths. In many countries, rates of premature births have increased between the 1990s and 2010s. Complications from preterm births resulted globally in 0.81 million deaths in 2015, down from 1.57 million in 1990. The chance of survival at 22 weeks is about 6%, while at 23 weeks it is 26%, 24 weeks 55% and 25 weeks about 72%. The chances of survival without any long-term difficulties are lower.
Signs and symptoms
Signs and symptoms of preterm labor include four or more uterine contractions in one hour. In contrast to false labour, true labor is accompanied by cervical dilation and effacement. Also, vaginal bleeding in the third trimester, heavy pressure in the pelvis, or abdominal or back pain could be indicators that a preterm birth is about to occur. A watery discharge from the vagina may indicate premature rupture of the membranes that surround the baby. While the rupture of the membranes may not be followed by labor, usually delivery is indicated as infection (chorioamnionitis) is a serious threat to both fetus and mother. In some cases, the cervix dilates prematurely without pain or perceived contractions, so that the mother may not have warning signs until very late in the birthing process.
Causes
The main categories of causes of preterm birth are preterm labor induction and spontaneous preterm labor.
Risk factors
The exact cause of spontaneous preterm birth is difficult to determine and it may be caused by many different factors at the same time as labor is a complex process. The research available is limited with regard to the cervix and therefore is limited in discerning what is or is not normal. Four different pathways have been identified that can result in preterm birth and have considerable evidence: precocious fetal endocrine activation, uterine overdistension (placental abruption), decidual bleeding, and intrauterine inflammation or infection.
Identifying women at high risk of giving birth early would enable the health services to provide specialized care for these women and their babies, for example a hospital with a special care baby unit such as a neonatal intensive care unit (NICU). In some instances, it may be possible to delay the birth. Risk scoring systems have been suggested as an approach to identify those at higher risk, however, there is no strong research in this area so it is unclear whether the use of risk scoring systems for identifying mothers would prolong pregnancy and reduce the numbers of preterm births or not.
Maternal factors
Risk factors in the mother have been identified that are linked to a higher risk of a preterm birth. These include age (either very young or older), high or low body mass index (BMI), length of time between pregnancies, endometriosis, previous spontaneous (i.e., miscarriage) or surgical abortions, unintended pregnancies, untreated or undiagnosed celiac disease, fertility difficulties, heat exposure, and genetic variables.
Studies on type of work and physical activity have given conflicting results, but it is opined that stressful conditions, hard labor, and long hours are probably linked to preterm birth. Obesity does not directly lead to preterm birth; however, it is associated with diabetes and hypertension which are risk factors by themselves. To some degree those individuals may have underlying conditions (i.e., uterine malformation, hypertension, diabetes) that persist. Couples who have tried more than one year versus those who have tried less than one year before achieving a spontaneous conception have an adjusted odds ratio of 1.35 (95% confidence interval 1.22–1.50) of preterm birth. Pregnancies after IVF confers a greater risk of preterm birth than spontaneous conceptions after more than one year of trying, with an adjusted odds ratio of 1.55 (95% CI 1.30–1.85).
Certain ethnicities may have a higher risk as well. For example, in the U.S. and the UK, Black women have preterm birth rates of 15–18%, more than double than that of the white population. Many Black women have higher preterm birth rates due to multiple factors but the most common is high amounts of chronic stress, which can eventually lead to premature birth. Adult chronic disease is not always the case with premature birth in Black women, which makes the main factor of premature birth challenging to identify. Filipinos are also at high risk of premature birth, and it is believed that nearly 11–15% of Filipinos born in the U.S. (compared to other Asians at 7.6% and whites at 7.8%) are premature. Filipinos being a big risk factor is evidenced with the Philippines being the eighth-highest ranking in the world for preterm births, the only non-African country in the top 10. This discrepancy is not seen in comparison to other Asian groups or Hispanic immigrants and remains unexplained. Genetic make-up is a factor in the causality of preterm birth. Genetics has been a big factor into why Filipinos have a high risk of premature birth as the Filipinos have a large prevalence of mutations that help them be predisposed to premature births. An intra- and transgenerational increase in the risk of preterm delivery has been demonstrated. No single gene has been identified.
Marital status has long been associated with risks for preterm birth. A 2005 study of 25,373 pregnancies in Finland revealed that unmarried mothers had more preterm deliveries than married mothers (P=0.001). Pregnancy outside of marriage was associated overall with a 20% increase in total adverse outcomes, even at a time when Finland provided free maternity care. A study in Quebec of 720,586 births from 1990 to 1997 revealed less risk of preterm birth for infants with legally married mothers compared with those with common-law wed or unwed parents. A study conducted in Malaysia in 2015 showed a similar trend, with marital status being significantly associated with preterm birth. However, the result of a study conducted in the US showed that between 1989 and 2006, marriage became less protective of preterm births which was attributed to the changing social norms and behaviors surrounding marriage.
Factors during pregnancy
Medications during pregnancy, living conditions, air pollution, smoking, illicit drugs or alcohol, infection, or physical trauma may also cause a preterm birth.
Air pollution: Living in an area with a high concentration of air pollution is a major risk factor for preterm labor, including living near major roadways or highways where vehicle emissions are high from traffic congestion or are a route for diesel trucks that tend to emit more pollution.
The use of fertility medication that stimulates the ovary to release multiple eggs and of IVF with embryo transfer of multiple embryos has been implicated as a risk factor for preterm birth. Often labor has to be induced for medical reasons; such conditions include high blood pressure, pre-eclampsia, maternal diabetes, asthma, thyroid disease, and heart disease.
Certain medical conditions in the pregnant mother may also increase the risk of preterm birth. Some women have anatomical problems that prevent the baby from being carried to term. These include a weak or short cervix (the strongest predictor of premature birth). Women with vaginal bleeding during pregnancy are at higher risk for preterm birth. While bleeding in the third trimester may be a sign of placenta previa or placental abruption—conditions that occur frequently preterm—even earlier bleeding that is not caused by these conditions is linked to a higher preterm birth rate. Women with abnormal amounts of amniotic fluid, whether too much (polyhydramnios) or too little (oligohydramnios), are also at risk. Anxiety and depression have been linked as risk factors for preterm birth.
The use of tobacco, cocaine, and excessive alcohol during pregnancy increases the chance of preterm delivery. Tobacco is the most commonly used drug during pregnancy and contributes significantly to low birth weight delivery. Babies with birth defects are at higher risk of being born preterm.
Passive smoking and/or smoking before the pregnancy influences the probability of a preterm birth. The World Health Organization published an international study in March 2014.
Presence of anti-thyroid antibodies is associated with an increased risk preterm birth with an odds ratio of 1.9 and 95% confidence interval of 1.1–3.5.
Intimate violence against the mother is another risk factor for preterm birth.
Physical trauma may case a preterm birth. The Nigerian cultural method of abdominal massage has been shown to result in 19% preterm birth among women in Nigeria, plus many other adverse outcomes for the mother and baby. This ought not be confused with massage therapy conducted by a fully trained and certified/licensed massage therapist or by significant others trained to provide massage during pregnancy, which—in a study involving pregnant females with prenatal depression—has been shown to have numerous positive results during pregnancy, including the reduction of preterm birth, less depression, lower cortisol, and reduced anxiety. In healthy women, however, no effects have been demonstrated in a controlled study.
Infection
The frequency of infection in preterm birth is inversely related to the gestational age. Mycoplasma genitalium infection is associated with increased risk of preterm birth, and spontaneous abortion.
Infectious microorganisms can be ascending, hematogenous, iatrogenic by a procedure, or retrograde through the fallopian tubes. From the deciduae they may reach the space between the amnion and chorion, the amniotic fluid, and the fetus. A chorioamnionitis also may lead to sepsis of the mother. Fetal infection is linked to preterm birth and to significant long-term disability including cerebral palsy.
It has been reported that asymptomatic colonization of the decidua occurs in up to 70% of women at term using a DNA probe suggesting that the presence of micro-organism alone may be insufficient to initiate the infectious response.
As the condition is more prevalent in black women in the U.S. and the UK, it has been suggested to be an explanation for the higher rate of preterm birth in these populations. It is opined that bacterial vaginosis before or during pregnancy may affect the decidual inflammatory response that leads to preterm birth. The condition known as aerobic vaginitis can be a serious risk factor for preterm labor; several previous studies failed to acknowledge the difference between aerobic vaginitis and bacterial vaginosis, which may explain some of the contradiction in the results.
Untreated yeast infections are associated with preterm birth.
A review into prophylactic antibiotics (given to prevent infection) in the second and third trimester of pregnancy (13–42 weeks of pregnancy) found a reduction in the number of preterm births in women with bacterial vaginosis. These antibiotics also reduced the number of waters breaking before labor in full-term pregnancies, reduced the risk of infection of the lining of the womb after delivery (endometritis), and rates of gonococcal infection. However, the women without bacterial vaginosis did not have any reduction in preterm births or pre-labor preterm waters breaking. Much of the research included in this review lost participants during follow-up so did not report the long-term effects of the antibiotics on mothers or babies. More research in this area is needed to find the full effects of giving antibiotics throughout the second and third trimesters of pregnancy.
A number of maternal bacterial infections are associated with preterm birth including pyelonephritis, asymptomatic bacteriuria, pneumonia, and appendicitis. A review into giving antibiotics in pregnancy for asymptomatic bacteriuria (urine infection with no symptoms) found the research was of very low quality but that it did suggest that taking antibiotics reduced the numbers of preterm births and babies with low birth weight. Another review found that one dose of antibiotics did not seem as effective as a course of antibiotics but fewer women reported side effects from one dose. This review recommended that more research is needed to discover the best way of treating asymptomatic bacteriuria.
A different review found that preterm births happened less for pregnant women who had routine testing for low genital tract infections than for women who only had testing when they showed symptoms of low genital tract infections. The women being routinely tested also gave birth to fewer babies with a low birth weight. Even though these results look promising, the review was only based on one study so more research is needed into routine screening for low genital tract infections.
Also periodontal disease has been shown repeatedly to be linked to preterm birth. In contrast, viral infections, unless accompanied by a significant febrile response, are considered not to be a major factor in relation to preterm birth.
Genetics
There is believed to be a maternal genetic component in preterm birth. Estimated heritability of timing-of-birth in women was 34%. However, the occurrence of preterm birth in families does not follow a clear inheritance pattern, thus supporting the idea that preterm birth is a non-Mendelian trait with a polygenic nature.
Prenatal care
The absence of prenatal care has been associated with higher rates of preterm births. Analysis of 15,627,407 live births in the United States in 1995–1998 concluded that the absence of prenatal care carried a 2.9 (95%CI 2.8, 3.0) times higher risk of preterm births. This same study found statistically significant relative risks of maternal anemia, intrapartum fever, unknown bleeding, renal disease, placental previa, hydramnios, placenta abruption, and pregnancy-induced hypertension with the absence of prenatal care. All these prenatal risks were controlled for other high-risk conditions, maternal age, gravidity, marital status, and maternal education. The absence of prenatal care prior to and during the pregnancy is primarily a function of socioeconomic factors (low family income and education), access to medical consultations (large distance from the place of residence to the healthcare unit and transportation costs), quality of healthcare, and social support. Efforts to decrease rates of preterm birth should aim to increase the deficits posed by the aforementioned barriers and to increase access to prenatal care.
Diagnosis
Placental alpha microglobulin-1
Placental alpha microglobulin-1 (PAMG-1) has been the subject of several investigations evaluating its ability to predict imminent spontaneous preterm birth in women with signs, symptoms, or complaints suggestive of preterm labor. In one investigation comparing this test to fetal fibronectin testing and cervical length measurement via transvaginal ultrasound, the test for PAMG-1 (commercially known as the PartoSure test) has been reported to be the single best predictor of imminent spontaneous delivery within 7 days of a patient presenting with signs, symptoms, or complaints of preterm labor. Specifically, the PPV, or positive predictive value, of the tests were 76%, 29%, and 30% for PAMG-1, fFN and CL, respectively (P < 0.01).
Fetal fibronectin
Fetal fibronectin (fFN) has become an important biomarker—the presence of this glycoprotein in the cervical or vaginal secretions indicates that the border between the chorion and decidua has been disrupted. A positive test indicates an increased risk of preterm birth, and a negative test has a high predictive value. It has been shown that only 1% of women in questionable cases of preterm labor delivered within the next week when the test was negative.
Ultrasound
Obstetric ultrasound has become useful in the assessment of the cervix in women at risk for premature delivery. A short cervix preterm is undesirable: A cervical length of less than at or before 24 weeks of gestational age is the most common definition of cervical incompetence.
Emerging Technologies
Technologies under research and development to facilitate earlier diagnosis of preterm births include sanitary pads that identify biomarkers such as fFN and PAMG-1 and others, when placed into the vagina. These devices then calculate a risk of preterm birth and send the findings to a smartphone. The notion that risk scoring systems are accurate in predicting preterm birth has been debated in multiple literature reviews.
Classification
In humans, the usual definition of preterm birth is birth before a gestational age of 37 complete weeks. In the normal human fetus, several organ systems mature between 34 and 37 weeks, and the fetus reaches adequate maturity by the end of this period. One of the main organs greatly affected by premature birth is the lungs. The lungs are one of the last organs to mature in the womb; because of this, many premature babies spend the first days and weeks of their lives on ventilators. Therefore, a significant overlap exists between preterm birth and prematurity. Generally, preterm babies are premature and term babies are mature. Preterm babies born near 37 weeks often have no problems relating to prematurity if their lungs have developed adequate surfactant, which allows the lungs to remain expanded between breaths. Sequelae of prematurity can be reduced to a small extent by using drugs to accelerate maturation of the fetus, and to a greater extent by preventing preterm birth.
Prevention
Historically efforts have been primarily aimed to improve survival and health of preterm infants (tertiary intervention). Such efforts, however, have not reduced the incidence of preterm birth. Increasingly primary interventions that are directed at all women, and secondary intervention that reduce existing risks are looked upon as measures that need to be developed and implemented to prevent the health problems of premature infants and children. Smoking bans are effective in decreasing preterm births. Different strategies are used in the administration of prenatal care, and future studies need to determine if the focus can be on screening for high-risk women, or widened support for low-risk women, or to what degree these approaches can be merged.
Before pregnancy
Adoption of specific professional policies can immediately reduce risk of preterm birth as the experience in assisted reproduction has shown when the number of embryos during embryo transfer was limited.
Many countries have established specific programs to protect pregnant women from hazardous or night-shift work and to provide them with time for prenatal visits and paid pregnancy-leave. The EUROPOP study showed that preterm birth is not related to type of employment, but to prolonged work (over 42 hours per week) or prolonged standing (over 6 hours per day). Also, night work has been linked to preterm birth. Health policies that take these findings into account can be expected to reduce the rate of preterm birth. Preconceptional intake of folic acid is recommended to reduce birth defects. There is also some evidence that folic acid supplement preconceptionally (before becoming pregnant) may reduce premature birth. Reducing smoking is expected to benefit pregnant women and their offspring.
During pregnancy
Self-care methods to reduce the risk of preterm birth include proper nutrition, avoiding stress, seeking appropriate medical care, avoiding infections, and the control of preterm birth risk factors (e.g. working long hours while standing on feet, carbon monoxide exposure, domestic abuse, and other factors). Reducing physical activity during pregnancy has not been shown to reduce the risk of a preterm birth. Healthy eating can be instituted at any stage of the pregnancy including nutritional adjustments and consuming suggested vitamin supplements. Calcium supplementation in women who have low dietary calcium may reduce the number of negative outcomes including preterm birth, pre-eclampsia, and maternal death. The World Health Organization (WHO) suggests 1.5–2 g of calcium supplements daily, for pregnant women who have low levels of calcium in their diet. Supplemental intake of C and E vitamins have not been found to reduce preterm birth rates.
While periodontal infection has been linked with preterm birth, randomized trials have not shown that periodontal care during pregnancy reduces preterm birth rates. Smoking cessation has also been shown to reduce the risk. The use of personal at home uterine monitoring devices to detect contractions and possible preterm births in women at higher risk of having a preterm baby have been suggested. These home monitors may not reduce the number of preterm births, however, using these devices may increase the number of unplanned antenatal visits and may reduce the number of babies admitted to special care when compared with women receiving normal antenatal care. Support from medical professionals, friends, and family during pregnancy may be beneficial at reducing caesarean birth and may reduce prenatal hospital admissions, however, these social supports alone may not prevent preterm birth.
Screening during pregnancy
Screening for asymptomatic bacteriuria followed by appropriate treatment reduces pyelonephritis and reduces the risk of preterm birth. Extensive studies have been carried out to determine if other forms of screening in low-risk women followed by appropriate intervention are beneficial, including screening for and treatment of Ureaplasma urealyticum, group B streptococcus, Trichomonas vaginalis, and bacterial vaginosis did not reduce the rate of preterm birth. Routine ultrasound examination of the length of the cervix may identify women at risk of preterm labour and tentative evidence suggests ultrasound measurement of the length of the cervix in those with preterm labor can help adjust management and results in the extension of pregnancy by about four days. Screening for the presence of fibronectin in vaginal secretions is not recommended at this time in women at low risk of preterm birth.
Reducing existing risks
Women are identified to be at increased risk for preterm birth on the basis of their past obstetrical history or the presence of known risk factors. Preconception intervention can be helpful in selected patients in a number of ways. Patients with certain uterine anomalies may have a surgical correction (i.e. removal of a uterine septum), and those with certain medical problems can be helped by optimizing medical therapies prior to conception, be it for asthma, diabetes, hypertension, and others.
Multiple pregnancies
In multiple pregnancies, which often result from use of assisted reproductive technology, there is a high risk of preterm birth. Selective reduction is used to reduce the number of fetuses to two or three.
Reducing indicated preterm birth
A number of agents have been studied for the secondary prevention of indicated preterm birth. Trials using low-dose aspirin, fish oil, vitamin C and E, and calcium to reduce preeclampsia demonstrated some reduction in preterm birth only when low-dose aspirin was used. Even if agents such as calcium or antioxidants were able to reduce preeclampsia, a resulting decrease in preterm birth was not observed.
Reducing spontaneous preterm birth
Reduction in activity by the mother—pelvic rest, limited work, bed rest—may be recommended although there is no evidence it is useful with some concerns it is harmful. Increasing medical care by more frequent visits and more education has not been shown to reduce preterm birth rates. Use of nutritional supplements such as omega-3 polyunsaturated fatty acids is based on the observation that populations who have a high intake of such agents are at low risk for preterm birth, presumably as these agents inhibit production of proinflammatory cytokines. A randomized trial showed a significant decline in preterm birth rates, and further studies are in the making.
Antibiotics
While antibiotics can get rid of bacterial vaginosis in pregnancy, this does not appear to change the risk of preterm birth. It has been suggested that chronic chorioamnionitis is not sufficiently treated by antibiotics alone (and therefore they cannot ameliorate the need for preterm delivery in this condition).
Progestogens
Progestogens—often given in the form of vaginal progesterone or hydroxyprogesterone caproate—relax the uterine musculature, maintain cervical length, and possess anti-inflammatory properties; all of which invoke physiological and anatomical changes considered to be beneficial in reducing preterm birth. Two meta-analyses demonstrated a reduction in the risk of preterm birth in women with recurrent preterm birth by 40–55%.
Progestogen supplementation also reduces the frequency of preterm birth in pregnancies where there is a short cervix. A short cervix is one that is less than 25mm, as detected during a transvaginal cervical length assessment in the midtrimester. However, progestogens are not effective in all populations, as a study involving twin gestations failed to see any benefit. Despite extensive research related to progestogen effectiveness, uncertainties remain concerning types of progesterone and routes of administration.
Cervical cerclage
In preparation for childbirth, the woman's cervix shortens. Preterm cervical shortening is linked to preterm birth and can be detected by ultrasonography. Cervical cerclage is a surgical intervention that places a suture around the cervix to prevent its shortening and widening. Numerous studies have been performed to assess the value of cervical cerclage and the procedure appears helpful primarily for women with a short cervix and a history of preterm birth. Instead of a prophylactic cerclage, women at risk can be monitored during pregnancy by sonography, and when shortening of the cervix is observed, the cerclage can be performed.
Treatment
Tertiary interventions are aimed at women who are about to go into preterm labor, or rupture the membranes or bleed preterm. The use of the fibronectin test and ultrasonography improves the diagnostic accuracy and reduces false-positive diagnosis. While treatments to arrest early labor where there is progressive cervical dilatation and effacement will not be effective to gain sufficient time to allow the fetus to grow and mature further, it may defer delivery sufficiently to allow the mother to be brought to a specialized center that is equipped and staffed to handle preterm deliveries. In a hospital setting women are hydrated via intravenous infusion (as dehydration can lead to premature uterine contractions).
If a baby has cardiac arrest at birth and is less than 22 to 24 weeks gestational age, attempts at resuscitation are not generally indicated.
Steroids
Severely premature infants may have underdeveloped lungs because they are not yet producing their own surfactant. This can lead directly to respiratory distress syndrome, also called hyaline membrane disease, in the neonate. To try to reduce the risk of this outcome, pregnant mothers with threatened premature delivery prior to 34 weeks are often administered at least one course of glucocorticoids, an antenatal steroid that crosses the placental barrier and stimulates the production of surfactant in the lungs of the baby. Steroid use up to 37 weeks is also recommended by the American Congress of Obstetricians and Gynecologists. Typical glucocorticoids that would be administered in this context are betamethasone or dexamethasone, often when the pregnancy has reached viability at 23 weeks.
In cases where premature birth is imminent, a second "rescue" course of steroids may be administered 12 to 24 hours before the anticipated birth. There are still some concerns about the efficacy and side effects of a second course of steroids, but the consequences of RDS are so severe that a second course is often viewed as worth the risk. A 2015 Cochrane review (updated in 2022) supports the use of repeat dose(s) of prenatal corticosteroids for women still at risk of preterm birth seven days or more after an initial course.
A Cochrane review from 2020 recommends the use of a single course of antenatal corticosteroids to accelerate fetal lung maturation in women at risk of preterm birth. Treatment with antenatal corticosteroids reduces the risk of perinatal death, neonatal death and respiratory distress syndrome and probably reduces the risk of IVH.
Concerns about adverse effects of prenatal corticosteroids include increased risk for maternal infection, difficulty with diabetic control, and possible long-term effects on neurodevelopmental outcomes for the infants. There is ongoing discussion about when steroids should be given (i.e. only antenatally or postnatally too) and for how long (i.e. single course or repeated administration). Despite these unknowns, there is a consensus that the benefits of a single course of prenatal glucocorticosteroids vastly outweigh the potential risks.
Antibiotics
The routine administration of antibiotics to all women with threatened preterm labor reduces the risk of the baby being infected with group B streptococcus and has been shown to reduce related mortality rates.
When membranes rupture prematurely, obstetrical management looks for development of labor and signs of infection. Prophylactic antibiotic administration has been shown to prolong pregnancy and reduced neonatal morbidity with rupture of membranes at less than 34 weeks. Because of concern about necrotizing enterocolitis, amoxicillin or erythromycin has been recommended but not amoxicillin + clavulanic acid (co-amoxiclav).
Tocolysis
A number of medications may be useful to delay delivery including: nonsteroidal anti-inflammatory drugs, calcium channel blockers, beta mimetics, and atosiban. Tocolysis rarely delays delivery beyond 24–48 hours. This delay, however, may be sufficient to allow the pregnant woman to be transferred to a center specialized for management of preterm deliveries and give administered corticosteroids to reduce neonatal organ immaturity. Meta-analyses indicate that calcium-channel blockers and an oxytocin antagonist can delay delivery by 2–7 days, and β2-agonist drugs delay by 48 hours but carry more side effects. Magnesium sulfate does not appear to be useful to prevent preterm birth. Its use before delivery, however, does appear to decrease the risk of cerebral palsy.
Mode of delivery
The routine use of caesarean section for early delivery of infants expected to have very low birth weight is controversial, and a decision concerning the route and time of delivery probably needs to be made on a case-by-case basis.
Neonatal care
In developed countries premature infants are usually cared for in a neonatal intensive care unit (NICU). The physicians who specialize in the care of very sick or premature babies are known as neonatologists. In the NICU, premature babies are kept under radiant warmers or in incubators (also called isolettes), which are bassinets enclosed in plastic with climate control equipment designed to keep them warm and limit their exposure to germs. Modern neonatal intensive care involves sophisticated measurement of temperature, respiration, cardiac function, oxygenation, and brain activity. After delivery, plastic wraps or warm mattresses are useful to keep the infant warm on their way to the NICU. Treatments may include fluids and nutrition through intravenous catheters, oxygen supplementation, mechanical ventilation support, and medications. In developing countries where advanced equipment and even electricity may not be available or reliable, simple measures such as kangaroo care (skin to skin warming), encouraging breastfeeding, and basic infection control measures can significantly reduce preterm morbidity and mortality. Kangaroo mother care (KMC) can decrease the risk of neonatal sepsis, hypothermia, hypoglycemia and increase exclusive breastfeeding. Bili lights may also be used to treat newborn jaundice (hyperbilirubinemia).
Water can be carefully provided to prevent dehydration but not so much to increase risks of side effects.
Breathing support
In terms of respiratory support, there may be little or no difference in the risk of death or chronic lung disease between high flow nasal cannulae (HFNC) and continuous positive airway pressure (CPAP) or nasal intermittent positive pressure ventilation (NPPV). For extremely preterm babies (born before 28 weeks' gestation), targeting a higher versus a lower oxygen saturation range makes little or no difference overall to the risk of death or major disability. Babies born before 32 weeks have been shown to have a lower risk of death from bronchopulmonary dysplasia if they have CPAP immediately after being born, compared to receiving either supportive care or assisted ventilation.
There is insufficient evidence for or against placing preterm stable twins in the same cot or incubator (co-bedding).
Nutrition
Meeting the appropriate nutritional needs of preterm infants is important for long-term health. Optimal care may require a balance of meeting nutritional needs and preventing complications related to feeding. The ideal growth rate is not known, however, preterm infants usually require a higher energy intake compared to babies who are born at term. The recommended amount of milk is often prescribed based on approximated nutritional requirements of a similar aged fetus who is not compromised. An immature gastrointestinal tract (GI tract), medical conditions (or co-morbidities), risk of aspirating milk, and necrotizing enterocolitis may lead to difficulties in meeting this high nutritional demand and many preterm infants have nutritional deficits that may result in growth restrictions. In addition, very small preterm infants cannot coordinate sucking, swallowing, and breathing. Tolerating a full enteral feeding (the prescribed volume of milk or formula) is a priority in neonatal care as this reduces the risks associated with venous catheters including infection, and may reduce the length of time the infant requires specialized care in the hospital. Different strategies can be used to optimize feeding for preterm infants. The type of milk/formula and fortifiers, route of administration (by mouth, tube feeding, venous catheter), timing of feeding, quantity of milk, continuous or intermittent feeding, and managing gastric residuals are all considered by the neonatal care team when optimizing care. The evidence in the form of high quality randomized trials is generally fairly weak in this area, and for this reason different neonatal intensive care units may have different practices and this results in a fairly large variation in practice. The care of preterm infants also varies in different countries and depends on resources that are available.
Human breast milk and formula
The American Academy of Pediatrics recommended feeding preterm infants human milk, finding "significant short- and long-term beneficial effects," including lower rates of necrotizing enterocolitis (NEC). In the absence of evidence from randomised controlled trials about the effects of feeding preterm infants with formula compared with mother's own breast milk, data collected from other types of studies suggest that mother's own breast milk is likely to have advantages over formula in terms of the baby's growth and development. Milk from human donors also reduces the risk of NEC by half in very low birth rate infants and very preterm infants.
Breast milk or formula alone may not be sufficient to meet the nutritional needs of some preterm infants. Fortification of breast milk or formula by adding extra nutrients is an approach often taken for feeding preterm infants, with the goal of meeting the high nutritional demand. High quality randomized controlled trials are needed in this field to determine the effectiveness of fortification. It is unclear if fortification of breast milk improves outcomes in preterm babies, though it may speed growth. Supplementing human milk with extra protein may increase short-term growth but the longer-term effects on body composition, growth and brain development are uncertain. Higher protein formula (between 3 and 4 grams of protein per kilo of body weight) may be more effective than low protein formula (less than 3 grams per kilo per day) for weight gain in formula-fed low-birth-weight infants. There is insufficient evidence about the effect on preterm babies' growth of supplementing human milk with carbohydrate, fat, and branched-chain amino acids. Conversely, there is some indication that preterm babies who cannot breastfeed may do better if they are fed only with diluted formula compared to full strength formula but the clinical trial evidence remains uncertain.
Individualizing the nutrients and quantities used to fortify enteral milk feeds in infants born with very low birth weight may lead to better short-term weight gain and growth but the evidence is uncertain for longer term outcomes and for the risk of serious illness and death. This includes targeted fortification (adjusting the level of nutrients in response to the results of a test on the breast milk) and adjustable fortification (adding nutrients based on testing the infant).
Multi-nutrient fortifier used to fortify human milk and formula has traditionally been derived from bovine milk. Fortifier derived from humans is available, however, the evidence from clinical trials is uncertain and it is not clear if there are any differences between human-derived fortifier and bovine-derived fortifier in terms of neonatal weight gain, feeding intolerance, infections, or the risk of death.
Timing of feeds
For very preterm infants, most neonatal care centres start milk feeds gradually, rather than starting with a full enteral feeding right away, however, is not clear if starting full enteral feeding early effects the risk of necrotising enterocolitis. In these cases, the preterm infant would be receiving the majority of their nutrition and fluids intravenously. The milk volume is usually gradually increased over the following weeks. Research into the ideal timing of enteral feeding and whether delaying enteral feeding or gradually introducing enteral feeds is beneficial at improving growth for preterm infants or low birth weight infants is needed. In addition, the ideal timing of enteral feeds to prevent side effects such as necrotising enterocolitis or mortality in preterm infants who require a packed red blood cell transfusion is not clear. Potential disadvantages of a more gradual approach to feeding preterm infants associated with less milk in the gut and include slower GI tract secretion of hormones and gut motility and slower microbial colonization of the gut.
Regarding the timing of starting fortified milk, preterm infants are often started on fortified milk/formula once they are fed 100 mL/kg of their body weight. Other some neonatal specialists feel that starting to feed a preterm infant fortified milk earlier is beneficial to improve intake of nutrients. The risks of feeding intolerance and necrotising enterocolitis related to early versus later fortification of human milk are not clear. Once the infant is able to go home from the hospital there is limited evidence to support prescribing a preterm (fortified) formula.
Intermittent feeding versus continuous feeding
For infants who weigh less than 1500 grams, tube feeding is usually necessary. Most often, neonatal specialists feed preterm babies intermittently with a prescribed amount of milk over a short period of time. For example, a feed could last 10–20 minutes and be given every 3 hours. This intermittent approach is meant to mimic conditions of normal bodily functions involved with feeding and allow for a cyclic pattern in the release of gastrointestinal tract hormones to promote development of the gastrointestinal system. In certain cases, continuous nasogastric feeding is sometimes preferred. There is low to very low certainty evidence to suggest that low birth weight babies who receive continuous nasogastic feeding may reach the benchmark of tolerating full enteral feeding later than babies fed intermittently and it is not clear if continuous feeding has any effect on weight gain or the number of interruptions in feedings. Continuous feeding may have little to no effect on length of body growth or head circumference and the effects of continuous feeding on the risk of developing necrotising enterocolitis is not clear.
Since preterm infants with gastro-oesophageal reflux disease do not have a fully developed antireflux mechanism, deciding on the most effective approach for nutrition is important. It is not clear if continuous bolus intragastric tube feeding is more effective compared to intermittent bolus intragastric tube feeding for feeding preterm infants with gastroesophageal reflux disease.
For infants who would benefit from intermittent bolus feeding, some infants may be fed using the "push feed" method using a syringe to gently push the milk or formula into the stomach of the infant. Others may be fed using a gravity feeding system where the syringe is attached directly to a tube and the milk or formula drips into the infant's stomach. It is not clear from medical studies which approach to intermittent bolus feeding is more effective or reduces adverse effects such as apnea, bradycardia, or oxygen desaturation episodes.
High volume feeds
High-volume (more than 180 mL per kilogram per day) enteral feeds of fortified or non-fortified human breast milk or formula may improve weight gain while the pre-term infant is hospitalized, however, there is insufficient evidence to determine if this approach improves growth of the neonate and other clinical outcomes including length of hospital stay. The risks or adverse effects associated with high-volume enteral feeding of preterm infants including aspiration pneumonia, reflux, apnea, and sudden oxygen desaturation episodes have not been reported in the trials considered in a 2021 systematic review.
Parenteral (intraveneous) nutrition
For preterm infants who are born after 34 weeks of gestation ("late preterm infants") who are critically ill and cannot tolerate milk, there is some weak evidence that the infant may benefit from including amino acids and fats in the intravenous nutrition at a later time point (72 hours or longer from hospital admission) versus early (less than 72 hours from admission to hospital), however further research is required to understand the ideal timing of starting intravenous nutrition.
Gastric residuals
For preterm infants in neonatal intensive care on gavage feeds, monitoring the volume and colour of gastric residuals, the milk and gastrointestinal secretions that remain in the stomach after a set amount of time, is common standard of care practice. Gastric residual often contains gastric acid, hormones, enzymes, and other substances that may help improve digestion and mobility of the gastrointestinal tract. Analysis of gastric residuals may help guide timing of feeds. Increased gastric residual may indicate feeding intolerance or it may be an early sign of necrotizing enterocolitis. Increased gastric residual may be caused by an underdeveloped gastrointestinal system that leads to slower gastric emptying or movement of the milk in the intestinal tract, reduced hormone or enzyme secretions from the gastrointestinal tract, duodenogastric reflux, formula, medications, and/or illness. The clinical decision to discard the gastric residuals (versus re-feeding) is often individualized based on the quantity and quality of the residual. Some experts also suggest replacing the fresh milk or curded milk and bile-stained aspirates, but not replacing haemorrhagic residual. Evidence to support or refute the practice of re-feeding preterm infants with gastric residuals is lacking.
Hyponatraemia and hypernatraemia
Imbalances of sodium (hyponatraemia and hypernatraemia) are common in babies born preterm. Hypernatraemia (sodium levels in the serum of more than 145-150 mmol/L) is common early on in preterm babies and the risk of hyponatraemia (sodium levels of less than 135 nmol/L) increases after about a week of birth if left untreated and prevention approaches are not used. Preventing complications associated with sodium imbalances is part of standard of care for preterm infants and includes careful monitoring of water and sodium given to the infant. The optimal sodium dose given immediately after birth (first day) is not clear and further research is needed to understand the idea management approach.
Hearing assessment
The Joint Committee on Infant Hearing (JCIH) state that for preterm infants who are in the neonatal intensive care unit (NICU) for a prolonged time should have a diagnostic audiologic evaluation before they are discharged from the hospital. Well babies follow a 1-2-3-month benchmark timeline where they are screened, diagnosed, and receiving intervention for a hearing loss. However, for very premature babies, it might not be possible to complete a hearing screen at one month of age due to several factors. Once the baby is stable, an audiologic evaluation should be performed. For premature babies in the NICU, auditory brainstem response (ABR) testing is recommended. If the infant does not pass the screen, they should be referred for an audiologic evaluation by an audiologist. If the infant is on aminoglycosides such as gentamicin for less than five days they should be monitored and have a follow-up 6–7 months of being discharged from the hospital to ensure there is no late onset hearing loss due to the medication.
Outcomes and prognosis
Preterm births can result in a range of problems including mortality and physical and mental delays.
Mortality and morbidity
In the U.S. where many neonatal infections and other causes of neonatal death have been markedly reduced, prematurity is the leading cause of neonatal mortality at 25%. Prematurely born infants are also at greater risk for having subsequent serious chronic health problems as discussed below.
The earliest gestational age at which the infant has at least a 50% chance of survival is referred to as the limit of viability. As NICU care has improved over the last 40 years, the limit of viability has reduced to approximately 24 weeks. Most newborns who die, and 40% of older infants who die, were born between 20 and 25.9 weeks (gestational age), during the second trimester.
As risk of brain damage and developmental delay is significant at that threshold even if the infant survives, there are ethical controversies over the aggressiveness of the care rendered to such infants. The limit of viability has also become a factor in the abortion debate.
Specific risks for the preterm neonate
Preterm infants usually show physical signs of prematurity in reverse proportion to the gestational age. As a result, they are at risk for numerous medical problems affecting different organ systems.
Neurological problems can include apnea of prematurity, hypoxic-ischemic encephalopathy (HIE), retinopathy of prematurity (ROP), developmental disability, transient hyperammonemia, cerebral palsy, and intraventricular hemorrhage, the latter affecting 25% of babies born preterm, usually before 32 weeks of pregnancy. Mild brain bleeds usually leave no or few lasting complications, but severe bleeds often result in brain damage or even death. Neurodevelopmental problems have been linked to lack of maternal thyroid hormones, at a time when their own thyroid is unable to meet postnatal needs.
Cardiovascular complications may arise from the failure of the ductus arteriosus to close after birth: patent ductus arteriosus (PDA).
Respiratory problems are common, specifically the respiratory distress syndrome (RDS or IRDS) (previously called hyaline membrane disease). Another problem can be chronic lung disease (previously called bronchopulmonary dysplasia or BPD).
Gastrointestinal and metabolic issues can arise from neonatal hypoglycemia, feeding difficulties, rickets of prematurity, hypocalcemia, inguinal hernia, and necrotizing enterocolitis (NEC).
Hematologic complications include anemia of prematurity, thrombocytopenia, and hyperbilirubinemia (jaundice) that can lead to kernicterus.
Infection, including sepsis, pneumonia, and urinary tract infection
Survival
The chance of survival at 22 weeks is about 6%, while at 23 weeks it is 26%, 24 weeks 55% and 25 weeks about 72% as of 2016. With extensive treatment up to 30% of those who survive birth at 22 weeks survive longer term as of 2019. The chances of survival without long-term difficulties is less. Of those who survive following birth at 22 weeks 33% have severe disabilities. In the developed world, overall survival is about 90% while in low-income countries survival rates are about 10%.
Some children will adjust well during childhood and adolescence, although disability is more likely nearer the limits of viability. A large study followed children born between 22 and 25 weeks until the age of 6 years old. Of these children, 46% had moderate to severe disabilities such as cerebral palsy, vision or hearing loss and learning disabilities, 34% had mild disabilities, and 20% had no disabilities; 12% had disabling cerebral palsy. Up to 15% of premature infants have significant hearing loss.
As survival has improved, the focus of interventions directed at the newborn has shifted to reduce long-term disabilities, particularly those related to brain injury. Some of the complications related to prematurity may not be apparent until years after the birth. A long-term study demonstrated that the risks of medical and social disabilities extend into adulthood and are higher with decreasing gestational age at birth and include cerebral palsy, intellectual disability, disorders of psychological development, behavior, and emotion, disabilities of vision and hearing, and epilepsy. Standard intelligence tests showed that 41% of children born between 22 and 25 weeks had moderate or severe learning disabilities when compared to the test scores of a group of similar classmates who were born at full term. It is also shown that higher levels of education were less likely to be obtained with decreasing gestational age at birth. People born prematurely may be more susceptible to developing depression as teenagers.
Some of these problems can be described as being within the executive domain and have been speculated to arise due to decreased myelinization of the frontal lobes. Studies of people born premature and investigated later with MRI brain imaging, demonstrate qualitative anomalies of brain structure and grey matter deficits within temporal lobe structures and the cerebellum that persist into adolescence. Throughout life they are more likely to require services provided by physical therapists, occupational therapists, or speech therapists. They are more likely to develop type 1 diabetes (roughly 1.2 times the rate) and type 2 diabetes (1.5 times).
Despite the neurosensory, mental and educational problems studied in school age and adolescent children born extremely preterm, the majority of preterm survivors born during the early years of neonatal intensive care are found to do well and to live fairly normal lives in young adulthood. Young adults born preterm seem to acknowledge that they have more health problems than their peers, yet feel the same degree of satisfaction with their quality of life.
Beyond the neurodevelopmental consequences of prematurity, infants born preterm have a greater risk for many other health problems. For instance, children born prematurely have an increased risk for developing chronic kidney disease.
Epidemiology
Preterm birth complicates 5–18% of births worldwide. In Europe and many developed countries the preterm birth rate is generally 5–9%, while in the U.S. from 2007 to 2022 the rate fluctuated from 9.6 to 10.5 per cent.
As weight is easier to determine than gestational age, the World Health Organization tracks rates of low birth weight (< 2,500 grams), which occurred in 16.5% of births in less developed regions in 2000. It is estimated that one third of these low birth weight deliveries are due to preterm delivery. Weight generally correlates to gestational age; however, infants may be underweight for other reasons than a preterm delivery. Neonates of low birth weight (LBW) have a birth weight of less than and are mostly but not exclusively preterm babies as they also include small for gestational age (SGA) babies. Weight-based classification further recognizes Very Low Birth Weight (VLBW) which is less than 1,500 g, and Extremely Low Birth Weight (ELBW) which is less than 1,000 g. Almost all neonates in these latter two groups are born preterm.
About 75% of nearly a million deaths due to preterm delivery would survive if provided warmth, breastfeeding, treatments for infection, and breathing support. Complications from preterm births resulted in 740,000 deaths in 2013, down from 1.57 million in 1990.
Society and culture
Economics
Preterm birth is a significant cost factor in healthcare, not even considering the expenses of long-term care for individuals with disabilities due to preterm birth. A 2003 study in the U.S. determined neonatal costs to be $224,400 for a newborn at 500–700 g versus $1,000 at over 3,000 g. The costs increase exponentially with decreasing gestational age and weight.
The 2007 Institute of Medicine report Preterm Birth found that the 550,000 premature babies born each year in the U.S. run up about $26 billion in annual costs, mostly related to care in neonatal intensive care units, but the real tab may top $50 billion.
Notable cases
James Elgin Gill (born on 20 May 1987 in Ottawa, Ontario, Canada) was the earliest premature baby in the world, until that record was broken in 2004. He was 128 days premature, 21 weeks 5 days gestation, and weighed . He survived.
In 2014, Lyla Stensrud, born in San Antonio, Texas, U.S., became the youngest premature baby in the world. She was born at 21 weeks 4 days and weighed 410 grams (less than a pound). Kaashif Ahmad resuscitated the baby after she was born. As of November 2018, Lyla was attending preschool. She had a slight delay in speech, but no other known medical issues or disabilities.
Amillia Taylor is also often cited as the most premature baby. She was born on 24 October 2006 in Miami, Florida, U.S., at 21 weeks and 6 days' gestation. This report has created some confusion as her gestation was measured from the date of conception (through in vitro fertilization) rather than the date of her mother's last menstrual period, making her appear 2 weeks younger than if gestation was calculated by the more common method. At birth, she was long and weighed . She had digestive and respiratory problems, together with a brain hemorrhage. She was discharged from the Baptist Children's Hospital on 20 February 2007.
The record for the smallest premature baby to survive was held for a considerable amount of time by Madeline Mann, who was born in 1989 at 26 weeks, weighing and measuring long. This record was broken in September 2004 by Rumaisa Rahman, who was born in the same hospital, Loyola University Medical Center in Maywood, Illinois. at 25 weeks' gestation. At birth, she was long and weighed . Her twin sister was also a small baby, weighing at birth. During pregnancy their mother had pre-eclampsia, requiring birth by caesarean section. The larger twin left the hospital at the end of December, while the smaller remained there until 10 February 2005 by which time her weight had increased to . Generally healthy, the twins had to undergo laser eye surgery to correct vision problems, a common occurrence among premature babies.
In May 2019, Sharp Mary Birch Hospital for Women & Newborns in San Diego announced that a baby nicknamed "Saybie" had been discharged almost five months after being born at 23 weeks' gestation and weighing . Saybie was confirmed by Dr. Edward Bell of the University of Iowa, which keeps the Tiniest Babies Registry, to be the new smallest surviving premature baby in that registry.
Born in February 2009, at Children's Hospitals and Clinics of Minnesota, Jonathon Whitehill was just 25 weeks' gestation with a weight of . He was hospitalized in a neonatal intensive care unit for five months, and then discharged.
Richard Hutchinson was born at Children's Hospitals and Clinics of Minnesota in Minneapolis, Minnesota, on June 5, 2020, at 21 weeks 2 days gestation. At birth he weighed . He remained hospitalized until November 2020, when he was then discharged.
On 5 July 2020 Curtis Means was born at the University of Alabama at Birmingham hospital at 21 weeks 1 day, and weighed . He was discharged in April 2021. , he is the current world record holder.
Historical figures who were born prematurely include Johannes Kepler (born in 1571 at seven months' gestation), Isaac Newton (born in 1642, small enough to fit into a quart mug, according to his mother), Winston Churchill (born in 1874 at seven months' gestation), and Anna Pavlova (born in 1885 at seven months' gestation).
Effect of the coronavirus pandemic
During the COVID-19 pandemic, a drastic drop in the rate of premature births has been reported in many countries, ranging from a 20% reduction to a 90% drop in the starkest cases. Studies in Ireland and Denmark first noticed the phenomenon, and it has been confirmed elsewhere. There is no universally accepted explanation for this drop as of August 2020. Hypotheses include additional rest and support for expectant mothers staying at home, less air pollution due to shutdowns and reduced car fumes, and reduced likelihood of catching other diseases and viruses in general due to the lockdowns.
Research
Brain injury is common among preterms, ranging from white matter injury to intraventricular and cerebellar haemorrhages. The characteristic neuropathology of preterms has been described as the "encephalopathy of prematurity". The number of preterms that receive special education is doubled compared to the general population. School marks are lower and so are verbal learning, executive function, language skills, and memory performance scores, as well as IQ scores. Behaviourally, adolescents who were born very preterm and/or very low birth weight have similar self-reports of quality of life, health status and self-esteem as term controls.
Various structural magnetic resonance studies found consistent reductions in whole brain volume. The extensive list of particular regions with smaller volumes compared to controls includes many cortical areas (temporal, frontal, parietal, occipital and cingulate), the hippocampal regions, thalamus, basal ganglia, amygdala, brain stem, internal capsule, corpus callosum, and cerebellum. Brain volume reduction seems to be present throughout the whole brain. In contrast, larger volumes were found in some of the same areas including medial/anterior frontal, parietal and temporal cortex, cerebellum, middle temporal gyrus, parahippocampal gyrus, and fusiform gyrus, as well as larger lateral ventricles on average. The cause of these inconsistencies are unknown. Additionally, reductions in cortical surface area/cortical thickness were found in the temporal lobes bilaterally and in left frontal and parietal areas. Thicker cortex was found bilaterally in the medial inferior and anterior parts of the frontal lobes and in the occipital lobes.
Gestational age was positively correlated with volumes of the temporal and fusiform gyri and sensorimotor cortex bilaterally, left inferior parietal lobule, brain stem, and various white matter tracts, as well as specific positive associations with the cerebellum and thalamus. Several structural brain alterations have been linked back to cognitive and behavioural outcome measures. For example, total brain tissue volume explained between 20 and 40% of the IQ and educational outcome differences between extremely preterm born adolescents and control adolescents. In another study, a 25% quartile decrease in white matter values in middle temporal gyrus was associated with a 60% increase in the risk of cognitive impairment. Nosarti and colleagues previously hypothesised that maturational patterns in preterm brains were consistent with the age-related stages typically observed in younger subjects. Their most recent study suggests, however, that their trajectory may not only be delayed but also fundamentally distinctive. Since both smaller and larger regional volumes were found in very preterm individuals compared to controls.
The evidence to support the use of osteopathic manipulations to provide benefit in neonatal care is weak.
| Biology and health sciences | Human reproduction | Biology |
429662 | https://en.wikipedia.org/wiki/Pentatomoidea | Pentatomoidea | The Pentatomoidea are a superfamily of insects in the suborder Heteroptera of the order Hemiptera. As hemipterans, they possess a common arrangement of sucking mouthparts. The roughly 7000 species under Pentatomoidea are divided into 21 families (16 extant and 5 extinct). Among these are the stink bugs and shield bugs, jewel bugs, giant shield bugs, and burrower bugs.
Description
The Pentatomoidea are characterised by a well-developed scutellum (the hardened extension of the thorax over the abdomen). It can be triangular to semielliptical in shape. The antennae typically have five segments. The tarsi usually have two or three segments.
Shield bugs have glands that produce a foul-smelling liquid, which is used defensively to deter potential predators. Nymphs have glands on the dorsal surface of the abdomen (dorsal abdominal scent glands). These are often present in adults as well, but adults also develop a pair of glands on the metathorax (third segment of the thorax), these being the metathoracic scent glands.
The nymphs and adults have distinctive piercing mouthparts, with mandibles and maxillae modified to form a piercing "stylet" sheathed within a modified labium. The stylet is used to suck sap from plants, or in some cases to suck blood from other animals, such as in the predatory subfamily Asopinae.
Pentatomoidea are mostly phytophagous, although some (the Asopinae or predatory stink bugs) are zoophagous. They can become significant pests (e.g. the brown marmorated stink bug), causing economic damage to certain crops.
Families
These families are classified under Pentatomoidea:
Extant
Acanthosomatidae – known as shield bugs, contains 46 genera and 184 species found worldwide
Canopidae – found strictly in the Neotropical realm
Cydnidae – known as burrowing bugs, it contains 120 genera and about 765 species worldwide.
Dinidoridae – found in tropical Asia, Africa, Australia, and South America, composed of 16 genera and about 65 species
Lestoniidae – small, round bugs that bear a resemblance to tortoise beetles (Chrysomelidae), composed only of one genus and two species, endemic to Australia
Megarididae – contains only one extant genus (Megaris) and 16 species, small, globular bugs occurring in Central America
Parastrachiidae – bright red and black bugs exhibiting maternal care of eggs, it contains only two genera: Dismegistus (Africa) and Parastrachia (Eastern Asia).
Pentatomidae – known as stink bugs, it is the largest family in Pentatomoidea. It contains around 900 genera and over 4700 species.
Phloeidae – large mottled brown and flattened bugs found strictly in the Neotropical realm. It is composed on only 2 genera and 3 species. They are known to exhibit strong maternal care.
Plataspidae – found in Asia, particularly eastern Asia, although a few species of Coptosoma occur in the Palearctic. They are round plant-feeding bugs. It has about 59 genera and 560 species.
Saileriolidae – only recently removed from inclusion within Urostylididae.
Scutelleridae – known as jewel bugs or shield-backed bugs. Composed of 81 genera and about 450 species.
Tessaratomidae – known as giant shield bugs because they are usually relatively large. Has about 55 genera and 240 species worldwide (mainly in the Old World tropics).
Thaumastellidae – small bugs usually found under rocks in tropical Africa and the Middle East. It contains only one genus and three species. There is some debate to their inclusion within Pentatomoidea.
Thyreocoridae – includes the former family, subfamily Corimelaeninae – known as ebony bugs, they are small, oval, shiny black bugs.
Urostylididae – contains about 11 genera and 170 species. They are found in Southern and Eastern Asia.
Extinct
†Mesopentacoridae Popov 1968 Middle Jurassic-Early Cretaceous, Asia
†Primipentatomidae – family with about four Early Cretaceous fossil species from China.
†Probascanionidae Handlirsch 1921 Monotypic, Early Jurassic, Germany
†Protocoridae Handlirsch 1906 Early-Middle Jurassic, Eurasia
†Venicoridae Yao et al. 2012 Early Cretaceous, China
Phylogeny
The morphological unweighted tree of Pentatomoidea after Grazia et al. (2008).
Gallery
| Biology and health sciences | Hemiptera (true bugs) | Animals |
429813 | https://en.wikipedia.org/wiki/Malacostraca | Malacostraca | Malacostraca is the second largest of the six classes of pancrustaceans behind insects, containing about 40,000 living species, divided among 16 orders. Its members, the malacostracans, display a great diversity of body forms and include crabs, lobsters, spiny lobsters, crayfish, shrimp, krill, prawns, woodlice, amphipods, mantis shrimp, tongue-eating lice and many other less familiar animals. They are abundant in all marine environments and have colonised freshwater and terrestrial habitats. They are segmented animals, united by a common body plan comprising 20 body segments (rarely 21), and divided into a head, thorax, and abdomen.
Etymology
The name Malacostraca is . The word was used by Aristotle, who contrasted them with oysters, in comparison with which their shells are pliable.
It was applied to this taxon by French zoologist Pierre André Latreille in 1802. He was curator of the arthropod collection at the National Museum of Natural History in Paris.
This scientific name is misleading, since the shell is soft only immediately after moulting, and is usually hard.
Malacostracans are sometimes contrasted with entomostracans, a name applied to all crustaceans outside the Malacostraca, and named after the obsolete taxon Entomostraca.
Description
The class Malacostraca includes about 40,000 species, and "arguably ... contains a greater diversity of body forms than any other class in the animal kingdom". Its members are characterised by the presence of three tagmata (specialized groupings of multiple segments) – a five-segmented head, an eight-segmented thorax and an abdomen with six segments and a telson, except in the Leptostraca, which retain the ancestral condition of seven abdominal segments. Malacostracans have abdominal appendages, a fact that differentiates them from all other major crustacean taxa except Remipedia. Each body segment bears a pair of jointed appendages, although these may be lost secondarily.
Tagmata
The head bears two pairs of antennae, the first of which is often biramous (branching into two parts) and the second pair bear exopods (outer branches) which are often flattened into antennal scales known as scaphocerites. The mouthparts consist of pairs each of mandibles, maxillules (second pair of mouthparts) and maxillae. Except for fairy shrimps, malacostracans are the only extant arthropods with compound eyes placed on moveable stalks, although in some taxa the eyes are unstalked, reduced or lost.
Up to three thoracic segments may be fused with the head to form a cephalothorax; the associated appendages turn forward and are modified as maxillipeds (accessory mouthparts). A carapace may be absent, present or secondarily lost, and may cover the head, part or all of the thorax and some of the abdomen. It is variable in form and may be fused dorsally with some of the thoracic segments or occasionally be in two parts, hinged dorsally. Typically, each of the thoracic appendages is biramous and the endopods are the better developed of the branches, being used for crawling or grasping. Each endopod consist of seven articulating segments; the coxa, basis, ischium, merus, carpus, propodus and dactylus. In decapods, the claw is formed by the articulation of the dactylus against an outgrowth of the propodus. In some taxa, the exopods are lost and the appendages are uniramous.
There is a clear demarcation between the thorax and the six or seven-segmented abdomen. In most taxa, each abdominal segment except the last carries a pair of biramous pleopods used for swimming, burrowing, gas exchange, creating a current or brooding eggs. The first and second abdominal pleopods may be modified in the male to form gonopods (accessory copulatory appendages). The appendages of the last segment are typically flattened into uropods, which together with the terminal telson, make up the "tail fan". It is the sudden flexion of this tail fan that provides the thrust for the rapid escape response of these crustaceans and the tail fan is also used in steering. In Leptostraca, the appendages on the telson instead form caudal rami (spine-like protrusions).
Internal anatomy
The digestive tract is straight and the foregut consists of a short oesophagus and a two-chambered stomach, the first part of which contains a gizzard-like "gastric mill" for grinding food. The walls of this have chitinous ridges, teeth and calcareous ossicles. The fine particles and soluble material are then moved into the midgut where chemical processing and absorption takes place in one or more pairs of large digestive caeca. The hindgut is concerned with water reclamation and the formation of faeces and the anus is situated at the base of the telson.
Like other crustaceans, malacostracans have an open circulatory system in which the heart pumps blood into the hemocoel (body cavity) where it supplies the needs of the organs for oxygen and nutrients before diffusing back to the heart. The typical respiratory pigment in malacostracans is haemocyanin. Structures that function as kidneys are located near the base of the antennae. A brain exists in the form of ganglia close to the antennae, there are ganglia in each segment and a collection of major ganglia below the oesophagus. Sensory organs include compound eyes (often stalked), ocelli (simple eyes), statocysts and sensory bristles. The naupliar eye is a characteristic of the nauplius larva and consists of four cup-shaped ocelli facing in different directions and able to distinguish between light and darkness.
Ecology
Malacostracans live in a wide range of marine and freshwater habitats, and three orders have terrestrial members: Amphipoda (Talitridae), Isopoda (Oniscidea, the woodlice) and Decapoda (terrestrial hermit crabs, crabs of the families Ocypodidae, Gecarcinidae, and Grapsidae, and terrestrial crayfish). They are abundant in all marine ecosystems, and most species are scavengers, although some, such as the porcelain crabs, are filter feeders, and some, such as mantis shrimps, are carnivores.
Life cycle
Most species of malacostracans have distinct sexes (a phenomenon known as gonochorism), although a few species exhibit hermaphroditism. The female genital openings or gonopores are located on the sixth thoracic segment or its appendages, while the male gonopores are on the eighth segment or its appendages, or in a small number of species, on the seventh. The naupliar larval stages are often reduced and take place before hatching, but where they occur, a metamorphosis usually occurs between the larval and the adult forms. Primitive malacostracans have a free-swimming naupliar larval stage. Research suggests the common ancestor of Malacostraca had lost the free-living nauplius larval stage, but re-evolved it again through heterochrony in Dendrobranchiata and Euphausiacea, which both has a lecithotrophic (non-feeding) nauplius stage.
Mating
Mating behavior has been studied in the freshwater shrimp Caridina ensifera. Multiple paternity, common in the Malacostrica, also occurs in C. ensifera. Reproductive success of sires was found to correlate inversely with their genetic relatedness to the mother. This finding suggests that sperm competition and/or pre- and post-copulatory female choice occurs. Female choice may increase the fitness of progeny by avoiding inbreeding that can lead to expression of homozygous deleterious recessive mutations.
Phylogenetics
The monophyly of Malacostraca is widely accepted. This is supported by several common morphological traits which are present throughout the group and is confirmed by molecular studies. However, a number of problems make it difficult to determine the relationships between the orders of Malacostraca. These include differences in mutation rates in different lineages, different patterns of evolution being apparent in different sources of data, including convergent evolution, and long branch attraction.
There is less agreement on the status of the subclass Phyllocarida with its single extant order, Leptostraca, depending on whether foliaceous (leaf-like) limbs have a single or multiple origin. Some authors advocate placing Phyllocarida in Phyllopoda, a group used in former classification systems, which would then include branchiopods, cephalocarids and leptostracans. A molecular study by American biologists Trisha Spears and Lawrence Abele concluded that phylogenetic evidence did not support the monophyly of this grouping, and that Phyllocarida should be regarded as a subclass of Malacostraca that had diverged from the main lineage at an early date.
The following cladogram is based on the 2001 phylogenetic analysis of Richter & Scholtz.
Subclass Phyllocarida
Leptostraca is the only extant order of Phyllocarida, the other two orders, Archaeostraca and Hoplostraca being extinct. Leptostracans are thought to be the most primitive of the malacostracans and date back to the Cambrian period. They range in length from , most being suspension feeders though some are carnivores or scavengers. They have a two part carapace which encloses the head, the whole thorax and part of the abdomen and are the only malacostracans with seven abdominal segments. Three families are known with several genera and about twenty species. They are found worldwide from the intertidal zone to the deep ocean, all but one species being benthic (living on the seabed).
Subclass Hoplocarida
Stomatopoda is the only extant order of Hoplocarida, the other two orders, Aeschronectida and Archaeostomatopoda being extinct. Stomatopodans, commonly known as mantis shrimps, range in length from and are predators. They have a dorso-ventrally flattened body and a shield-like carapace and are armed with powerful, raptorial claws normally carried in a folded position. There are about 300 species, most living in tropical and subtropical seas although some live in temperate areas. They are benthic, mostly hiding in cracks and crevices or living in burrows, some emerging to forage while others are ambush predators.
Subclass Eumalacostraca
The Eumalocostraca contains the vast majority of the approximately 40,000 living species of malacostracans and consists of three superorders, Syncarida, Peracarida and Eucarida. Syncaridans are mostly small and found in freshwater and subterranean habitats. Peracaridans are characterised by having a marsupium in which they brood their young. They are found in marine, freshwater and terrestrial habitats and include Amphipoda, Cumacea, Isopoda and Mysida. Eucarida includes lobsters, crabs, shrimps, prawns and krill.
Fossil record
The first malacostracans appeared sometime in the Cambrian, when animals belonging to the Phyllocarida appeared.
Classification
The following classification of living malacostracans is based on An Updated Classification of the Recent Crustacea (2001) by the American marine biologists Joel W. Martin, curator of crustaceans at the Natural History Museum of Los Angeles County, and George E. Davies. Extinct orders have been added to this and are indicated by an obelisk (†).
Class Malacostraca Latreille, 1802
Subclass Phyllocarida Packard, 1879
† Archaeostraca Claus 1888
† Hoplostraca Schram, 1973
Leptostraca Claus, 1880
Subclass Hoplocarida Calman, 1904
† Aeschronectida Schram, 1969
† Archaeostomatopoda Schram, 1969
Stomatopoda Latreille, 1817
Subclass Eumalacostraca Grobben, 1892
Superorder Syncarida Packard, 1885
† Palaeocaridacea Brooks, 1979
Bathynellacea Chappuis, 1915
Anaspidacea Calman, 1904
Superorder Peracarida Calman, 1904
Spelaeogriphacea Gordon, 1957
Thermosbaenacea Monod, 1927
Lophogastrida Sars, 1870
Mysida Haworth, 1825
Stygiomysida Tchindonova, 1981
Mictacea Bowman et al., 1985
Amphipoda Latreille, 1816
Isopoda Latreille, 1817
Tanaidacea Dana, 1849
Cumacea Krøyer, 1846
Superorder Eucarida Calman, 1904
† Angustidontida Gueriau, Charbonnier & Clément, 2014
Euphausiacea Dana, 1852
Decapoda Latreille, 1802
| Biology and health sciences | Crustaceans | null |
429917 | https://en.wikipedia.org/wiki/Angora%20rabbit | Angora rabbit | The Angora rabbit (), one of the most ancient groups of domestic rabbit breeds, which is bred for the long fibers of its coat, known as Angora wool. They are gathered by shearing, combing or plucking. Because rabbits do not possess the same allergy-causing qualities as many other animals, their wool is an important alternative. There are at least 11 distinct breeds of Angora rabbit, four of which are currently recognized by the American Rabbit Breeders Association (ARBA): the English Angora, the French Angora, the Giant Angora and the Satin Angora. Other unrecognized breeds include the German Angora, the Finnish Angora, the Chinese Angora, the Japanese Angora, the Korean Angora, the Russian Angora, the St Lucian Angora and the Swiss Angora.
History
The Angora is said to have originated in Ankara (historically known as Angora), in present-day Turkey, and is known to have been brought to France in 1723. The Angora rabbit became a popular pet of the French royalty in the mid-18th century, and Angoras had spread to other parts of Europe by the end of that century. In the United States, garments made of Angora-rabbit wool have been popular ever since they first arrived in the early 20th century. However, only during World War II did domestic production expand to meet the demand for more than a year. This valuable, soft, silky, fiber aroused much interest, and quickly people became enamored with the production process.
Angora-rabbit wool
Angoras are bred mainly for their wool, which is silky and soft. At only 14–16 micrometres in diameter, it is similar to cashmere in fineness and softness to the touch. A healthy adult Angora's wool will grow approximately per month. Regular grooming is necessary to prevent the fibre from matting and felting on the rabbit, which causes discomfort and can lead to pain and even infection. Angora wool is harvested (plucked or shorn) every three to four months throughout the year. The coat needs to be monitored after 6 months of regrowth since it may tend to "die" and easily mat.
Angora wool may be gathered periodically by hand-plucking the hair within the coat that are being naturally shed. Full harvesting is done by shearing the coat with clippers or small scissors, often while the rabbit sits atop a groomer's turntable. Shearing typically starts at the head, moving across the shoulders to the tail. The rabbit is then flipped and the underside is shorn from tail to chin. Between of wool may be harvested from a Giant Angora.
Health
Wool block
Because of the length and abundance of their hair, Angora rabbits are particularly susceptible to wool block, a potentially lethal blockage of the digestive tract. All rabbits ingest some of their wool when they groom themselves, but their digestive system is sometimes not able to pass that foreign matter. The length of Angora hair compounds the risk of impaction, which can lead to death. Clipping or plucking an Angora's wool every 90–120 days is necessary to prevent wool block.
Wool mites
Cheyletiella parasitovorax is a skin parasite commonly found in Angora rabbits. Signs of infestation are flaky skin patches and fur loss. Wool mites reduce fiber yields and the resulting skin flakes are detrimental to the fiber quality. Wool mites may be treated with ivermectin or with carbaryl powder.
Angora rabbit breeds
The iconic long coat of the Angora is the result of a rabbit gene referred to as l (i.e., lowercase "L"). This "Angora gene" is present in all Angora breeds. It has also sometimes been utilized in the development of other rabbit breeds or other breeds' new varieties. "Dwarf Wooly" breeds including American Fuzzy Lop, Lionhead and Jersey Wooly are now recognized in the U.S. by ARBA. Belgium and France have their own Dwarf Wooly breeds. There is also a rare Mini English Angora breed in New Zealand.
English Angora
Weight:
ARBA-recognized varieties: Agouti, Broken, Pointed White, Ruby-eyed White, Self, and Shaded
Before 1939 there was one breed of "Angora Wooler". In 1939 ARBA reclassified 'Angora Wooler' as English Type and French Type. In 1944 ARBA officially separated Angora rabbits into two breeds: English Angora and French Angora.
Rabbits of the English Angora breed are adorned with "fur", growths of wool on the ears and the entire face except above the nose, and front feet, along with their thick body, and wool. They are gentle in nature, but they are not recommended for those who do not groom their animals. Their wool is very dense and needs to be groomed twice a week.
This is the smallest Angora rabbit of the four ARBA-recognized breeds. This breed is more common as a pet because of the facial features that give it a puppy or teddy-bear look. If the texture of the wool is correct, the maintenance is relatively easy; if the texture of the rabbit is cottony, it requires a great deal of maintenance. Beginning spinners may find the wool challenging.
The English Angora can be bred to have broken colors—i.e., white with black spots—but this is not accepted by ARBA standards and would lead to a disqualification when showing the rabbit. When an English Angora rabbit is shown, the toenails should also be only one color, the ears could be folded over at the tips and the furnishings on the face may cover their eyes. The English Angora is the only one of the Angora breeds that has hair covering its eyes.
French Angora
Weight:
ARBA-recognized varieties: Agouti, Broken, Pointed White, Self, Shaded, Ticked, and Wide Band
This breed has a dense undercoat. If the texture is correct it requires less maintenance than other Angora breeds. Small ear tufts are allowed but not usually preferred by breeders. ARBA recognizes the same colors as with English Angora, plus ticked and wide band. They are shown at ARBA shows using the types 'white' and 'colored' (broken being a colored). As with other ARBA-shown rabbits, toenails should also be only one color.
The French Angora is one of the large Angora breeds at , with a commercial body type. It differs from the English, Giant and German Angora in that it possesses a clean (hairless) face and front feet with only minor tufting on the rear legs. The color of a French Angora is determined by the color of its head, feet and tail (all the same color). This variety of angora fibre has smooth silky texture. Beginning spinners may find Angora wool a challenge.
Desirable characteristics of the fibre include its texture, warmth, light weight and pure white color. It is used for sweaters, mittens, baby clothes and millinery.
German Angora
Weight:
IAGARB-accepted varieties: Albino or Colored (but not bi-colored)
Albino
Black
Dilute Blacka/k/a Blue
Browna/k/a Chocolate
Dilute Browna/k/a Lilac
Tortoiseshell
Dilute Black Tortoiseshella/k/a Blue Tortoiseshell
Brown Tortoiseshella/k/a Chocolate Tortoiseshell
Dilute Brown Tortoiseshella/k/a Lilac Tortoiseshell
Agoutia/k/a Black Agouti, Chestnut Agouti, Wild Agouti
Dilute Black Agoutia/k/a Blue Agouti or Opal
Brown Agoutia/k/a Chocolate Agouti
Dilute Brown Agoutia/k/a Lilac Agouti or Lynx
Yellowa/k/a Red or Fawn
Chinchilla
Dilute Black Chinchillaa/k/a Blue Chinchilla or Squirrel
Brown Chinchillaa/k/a Chocolate Chinchilla
Dilute Brown Chinchillaa/k/a Lilac Chinchilla
Though common, the German Angora is not currently recognized by ARBA. The International Association of German Angora Rabbit Breeders (IAGARB) maintains a breed standard for the German Angora.
Giant Angora
Weight: Minimum
ARBA-recognized varieties: Ruby-eyed White
The Giant Angora is the largest of the ARBA-recognized Angora breeds. It was originally developed to be an efficient commercial producer that could be sustained on 16–18% protein pellets plus hay, and live in the standard sized, all-wire cages.
Because ARBA wouldn't allow German Angoras to be shown, their body type being considered too similar to other Angora breeds, Louise Walsh of Taunton, Massachusetts, created a new breed. She used German Angoras, French Lops and Flemish Giants to develop a completely different 'commercial' body type. ARBA officially recognized the Giant Angora in 1988. Its coat includes three types of wool: soft underwool, awn fluff and awn hair.
The awn-type wool exists only in the Giant and German Angora breeds. The Giant Angora has furnishings on the face and ears. Many people confuse the German with the Giant Angora, but it is their body type that differs.
The only color variety ARBA currently recognizes for the Giant Angora is the Ruby-eyed White (REW), a color that indicates the genetic absence of pigment (albino). The Giant Angora produces more wool than the French, Satin or English Angoras. The Giant Angora is the only 6-class animal in the Angora breed. It should have a commercial-type body with a very dense coat of wool. The head should be oval in appearance, that is broad across the forehead and slightly narrower at the muzzle. The Giant Angora should have forehead tufts (head trimmings) and cheek furnishings. The head trimmings should be noticeable, but does have lighter trimmings than bucks. The ears should be lightly fringed and well tasseled. The Giant Angora is also the only breed of angora that is shown only as a ruby-eyed white. A Black color variety of the Giant Angora is in development but has not been sanctioned by ARBA.
The Giant Angora coat contains three fiber types for its texture. The underwool should be the most dominant over the other two types of hair. It should be medium-fine, soft and delicately waved and have a gentle shine. Beginning spinners may find Angora wool a challenge.
The Awn Fluff has a guard hair tip and is a stronger, wavy wool. The Awn Fluff is found between the underwool and Awn Hair. The Awn Hair, also known as guard hair, is the third type of fiber. The Awn Hair is a strong straight hair that protrudes above the wool and must be present and evident.
The classification of the Giant Angora is different from the other three breeds owing to it being a 6-class animal. The junior buck and junior doe must be under 6 months of age and have a minimum weight of . The intermediate buck and intermediate doe are 6–8 months of age. The senior buck and senior doe are 8 months of age or over. The senior buck must weigh at least . The senior doe must weigh at least .
When Giant Angoras are judged the majority of the points are based on the wool, which includes density, texture and length. The points for 'general type' include the body type, head, ears, eyes, feet, legs and tail.
Like many other 'giant' breeds of rabbit, the Giant Angora grows slowly. A doe usually takes more than a year to reach maturity (size and weight). A buck can take up to 1.5 years to mature (size and weight).
Satin Angora
Weight:
ARBA-recognized varieties: [Includes eight color groups. The color of a Satin Angora is determined by the uniform pigment on its head, feet, and tail.]
Pointed White (includes Black, Blue, Chocolate, Lilac)
Blue-Eyed White
Red-Eyed White
Broken
Chestnut
Chocolate Agouti
Chocolate Chinchilla
Chinchilla
Copper
Lilac Chinchilla
Lynx
Opal
Squirrel
Black
Blue
Chocolate
Lilac
Pearl
Black Pearl
Blue Pearl
Chocolate Pearl
Lilac Pearl
Sable Pearl
Smoke Pearl
Sable
Seal
Blue Tortoiseshell
Chocolate Tortoiseshell
Lilac Tortoiseshell
Tortoiseshell
Blue Steel
Chocolate Steel
Lilac Steel
Steel
Cream
Fawn
Red
The Satin Angora was developed in the late 1970s by Mrs. Meyer of Holland Landing, Ontario, Canada, who crossed French Angoras with rabbits of the Satin breed. In addition to the sheen (for which the Satin is known), true red and copper pigments emerged in the new rabbits. In all 'satinized' coats, the hair shaft has a semi-transparent outer shell that reflects light, resulting in deep color, high luster and an extremely soft and silky texture to the hair.
The Satin Angora (like the French Angora) has no furnishings on the face, ears or feet. The Satin does not produce as much wool as other Angora breeds, but this trait is being improved upon through selective breeding. While more difficult to keep groomed than the Giant or French Angora, the Satin is less difficult than the English Angora. Because of the soft texture of the wool and the lower guard-hair count in the coat, matting occurs more readily. Daily combing is therefore recommended.
Satin Angora wool is said to be stronger for spinning than other Angora varieties, but because of its slipperiness it can be more difficult to spin.
Other Angora rabbit breeds
Finnish Angora
Japanese Angora
Russian Angora
St. Lucian Angora
Swiss Angora
Dutch Angora
Genetics
Some genes or rather mutations causing the Angora phenotype have been identified. A gene that has been repeatedly found to be affected in Angora rabbits is the FGF5 gene. For example, a specific mutation (T19234C) changes the amino acid threonine (T) to cysteine (C) in the Fgf5 protein, causing the phenotype.
| Biology and health sciences | Rabbits | Animals |
430014 | https://en.wikipedia.org/wiki/Fused%20quartz | Fused quartz | Fused quartz, fused silica or quartz glass is a glass consisting of almost pure silica (silicon dioxide, SiO2) in amorphous (non-crystalline) form. This differs from all other commercial glasses, such as soda-lime glass, lead glass, or borosilicate glass, in which other ingredients are added which change the glasses' optical and physical properties, such as lowering the melt temperature, the spectral transmission range, or the mechanical strength. Fused quartz, therefore, has high working and melting temperatures, making it difficult to form and less desirable for most common applications, but is much stronger, more chemically resistant, and exhibits lower thermal expansion, making it more suitable for many specialized uses such as lighting and scientific applications.
The terms fused quartz and fused silica are used interchangeably but can refer to different manufacturing techniques, resulting in different trace impurities. However fused quartz, being in the glassy state, has quite different physical properties compared to crystalline quartz despite being made of the same substance. Due to its physical properties it finds specialty uses in semiconductor fabrication and laboratory equipment, for instance.
Compared to other common glasses, the optical transmission of pure silica extends well into the ultraviolet and infrared wavelengths, so is used to make lenses and other optics for these wavelengths. Depending on manufacturing processes, impurities will restrict the optical transmission, resulting in commercial grades of fused quartz optimized for use in the infrared, or in the ultraviolet. The low coefficient of thermal expansion of fused quartz makes it a useful material for precision mirror substrates or optical flats.
Manufacture
Fused quartz is produced by fusing (melting) high-purity silica sand, which consists of quartz crystals. There are four basic types of commercial silica glass:
Type I is produced by induction melting natural quartz in a vacuum or an inert atmosphere.
Type II is produced by fusing quartz crystal powder in a high-temperature flame.
Type III is produced by burning SiCl4 in a hydrogen-oxygen flame.
Type IV is produced by burning SiCl4 in a water vapor-free plasma flame.
Quartz contains only silicon and oxygen, although commercial quartz glass often contains impurities. Two dominant impurities are aluminium and titanium which affect the optical transmission at ultraviolet wavelengths. If water is present in the manufacturing process, hydroxyl (OH) groups may become embedded which reduces transmission in the infrared.
Fusion
Melting is effected at approximately 2200 °C (4000 °F) using either an electrically heated furnace (electrically fused) or a gas/oxygen-fuelled furnace (flame-fused). Fused silica can be made from almost any silicon-rich chemical precursor, usually using a continuous process which involves flame oxidation of volatile silicon compounds to silicon dioxide, and thermal fusion of the resulting dust (although alternative processes are used). This results in a transparent glass with an ultra-high purity and improved optical transmission in the deep ultraviolet. One common method involves adding silicon tetrachloride to a hydrogen–oxygen flame.
Product quality
Fused quartz is normally transparent. The material can, however, become translucent if small air bubbles are allowed to be trapped within. The water content (and therefore infrared transmission) of fused quartz is determined by the manufacturing process. Flame-fused material always has a higher water content due to the combination of the hydrocarbons and oxygen fueling the furnace, forming hydroxyl [OH] groups within the material. An IR grade material typically has an [OH] content below 10 ppm.
Applications
Many optical applications of fused quartz exploit its wide transparency range, which can extend well into the ultraviolet and into the near-mid infrared. Fused quartz is the key starting material for optical fiber, used for telecommunications.
Because of its strength and high melting point (compared to ordinary glass), fused quartz is used as an envelope for halogen lamps and high-intensity discharge lamps, which must operate at a high envelope temperature to achieve their combination of high brightness and long life. Some high-power vacuum tubes used silica envelopes whose good transmission at infrared wavelengths facilitated radiation cooling of their incandescent anodes.
Because of its physical strength, fused quartz was used in deep diving vessels such as the bathysphere and benthoscope and in the windows of crewed spacecraft, including the Space Shuttle and International Space Station. Fused quartz was used also in composite armour development.
In the semiconductor industry, its combination of strength, thermal stability, and UV transparency makes it an excellent substrate for projection masks for photolithography. Its UV transparency also finds use as windows on EPROMs (erasable programmable read only memory), a type of non-volatile memory chip which is erased by exposure to strong ultraviolet light. EPROMs are recognizable by the transparent fused quartz (although some later models use UV-transparent resin) window which sits on top of the package, through which the silicon chip is visible, and which transmits UV light for erasing.
Due to the thermal stability and composition, it is used in 5D optical data storage and in semiconductor fabrication furnaces.
Fused quartz has nearly ideal properties for fabricating first surface mirrors such as those used in telescopes. The material behaves in a predictable way and allows the optical fabricator to put a very smooth polish onto the surface and produce the desired figure with fewer testing iterations. In some instances, a high-purity UV grade of fused quartz has been used to make several of the individual uncoated lens elements of special-purpose lenses including the Zeiss 105 mm f/4.3 UV Sonnar, a lens formerly made for the Hasselblad camera, and the Nikon UV-Nikkor 105 mm f/4.5 (presently sold as the Nikon PF10545MF-UV) lens. These lenses are used for UV photography, as the quartz glass can be transparent at much shorter wavelengths than lenses made with more common flint or crown glass formulas.
Fused quartz can be metallised and etched for use as a substrate for high-precision microwave circuits, the thermal stability making it a good choice for narrowband filters and similar demanding applications. The lower dielectric constant than alumina allows higher impedance tracks or thinner substrates.
Refractory material applications
Fused quartz as an industrial raw material is used to make various refractory shapes such as crucibles, trays, shrouds, and rollers for many high-temperature thermal processes including steelmaking, investment casting, and glass manufacture. Refractory shapes made from fused quartz have excellent thermal shock resistance and are chemically inert to most elements and compounds, including virtually all acids, regardless of concentration, except hydrofluoric acid, which is very reactive even in fairly low concentrations. Translucent fused-quartz tubes are commonly used to sheathe electric elements in room heaters, industrial furnaces, and other similar applications.
Owing to its low mechanical damping at ordinary temperatures, it is used for high-Q resonators, in particular, for wine-glass resonator of hemispherical resonator gyro. For the same reason fused quartz is also the material used for modern glass instruments such as the glass harp and the verrophone, and is also used for new builds of the historical glass harmonica, giving these instruments a greater dynamic range and a clearer sound than with the historically used lead crystal.
Quartz glassware is occasionally used in chemistry laboratories when standard borosilicate glass cannot withstand high temperatures or when high UV transmission is required. The cost of production is significantly higher, limiting its use; it is usually found as a single basic element, such as a tube in a furnace, or as a flask, the elements in direct exposure to the heat.
Properties of fused quartz
The extremely low coefficient of thermal expansion, about (20–320 °C), accounts for its remarkable ability to undergo large, rapid temperature changes without cracking (see thermal shock).
Fused quartz is prone to phosphorescence and "solarisation" (purplish discoloration) under intense UV illumination, as is often seen in flashtubes. "UV grade" synthetic fused silica (sold under various tradenames including "HPFS", "Spectrosil", and "Suprasil") has a very low metallic impurity content making it transparent deeper into the ultraviolet. An optic with a thickness of 1 cm has a transmittance around 50% at a wavelength of 170 nm, which drops to only a few percent at 160 nm. However, its infrared transmission is limited by strong water absorptions at 2.2 μm and 2.7 μm.
"Infrared grade" fused quartz (tradenames "Infrasil", "Vitreosil IR", and others), which is electrically fused, has a greater presence of metallic impurities, limiting its UV transmittance wavelength to around 250 nm, but a much lower water content, leading to excellent infrared transmission up to 3.6 μm wavelength. All grades of transparent fused quartz/fused silica have nearly identical mechanical properties.
Refractive index
The optical dispersion of fused quartz can be approximated by the following Sellmeier equation:
where the wavelength is measured in micrometers. This equation is valid between 0.21 and 3.71 μm and at 20 °C. Its validity was confirmed for wavelengths up to 6.7 μm. Experimental data for the real (refractive index) and imaginary (absorption index) parts of the complex refractive index of fused quartz reported in the literature over the spectral range from 30 nm to 1000 μm have been reviewed by Kitamura et al. and are available online.
Its quite high Abbe Number of 67.8 makes it among the lowest dispersion glasses at visible wavelengths, as well as having an exceptionally low refractive index in the visible (nd = 1.4585). Note that fused quartz has a very different and lower refractive index compared to crystalline quartz which is birefringent with refractive indices no = 1.5443 and ne = 1.5534 at the same wavelength. Although these forms have the same chemical formula, their differing structures result in different optical and other physical properties.
List of physical properties
Density: 2.203 g/cm3
Hardness: 5.3–6.5 (Mohs scale), 8.8 GPa
Tensile strength: 48.3 MPa
Compressive strength: > 1.1 GPa
Bulk modulus: ~37 GPa
Rigidity modulus: 31 GPa
Young's modulus: 71.7 GPa
Poisson's ratio: 0.17
Lamé elastic constants: λ = 15.87 GPa, μ = 31.26 GPa
Coefficient of thermal expansion: 5.5 × 10−7/K (average 20–320 °C)
Thermal conductivity: 1.3 W/(m·K)
Specific heat capacity: 45.3 J/(mol·K)
Softening point: ≈ 1665 °C
Annealing point: ≈ 1140 °C
Strain point: 1070 °C
Electrical resistivity: > 1018 Ω·m
Dielectric constant: 3.75 at 20 °C 1 MHz
Dielectric loss factor: less than 0.0004 at 20 °C 1 MHz typically 6 × 10−5 at 10 GHz
Dielectric strength: 250–400 kV/cm at 20 °C
Magnetic susceptibility: −11.28 × 10−6 (SI, 22 °C)
Hamaker constant: A = 6.5 × 10−20 J.
Surface tension: 0.300 N/m at 1800–2400 °C
Index of refraction: nd = 1.4585 (at 587.6 nm)
Change of refractive index with temperature: 1.28 × 10−5/K (20–30 °C)
Transmission range: Cutoff – 160 to 5000 nm, with a deep absorption band at 2730 nm. Best transmittance – 180 to 2700 nm.
Stress-optic coefficients: p11 = 0.113, p12 = 0.252.
Abbe number: Vd = 67.82
| Technology | Materials | null |
430306 | https://en.wikipedia.org/wiki/Vampire%20bat | Vampire bat | Vampire bats, members of the subfamily Desmodontinae, are leaf-nosed bats currently found in Central and South America. Their food source is the blood of other animals, a dietary trait called hematophagy. Three extant bat species feed solely on blood: the common vampire bat (Desmodus rotundus), the hairy-legged vampire bat (Diphylla ecaudata), and the white-winged vampire bat (Diaemus youngi). Two extinct species of the genus Desmodus have been found in North America.
Taxonomy
Due to differences among the three species, each has been placed within a different genus, each consisting of one extant species. In the older literature, these three genera were placed within a family of their own, Desmodontidae, but taxonomists have now grouped them as a subfamily, Desmodontinae, in the New World leaf-nosed bat family, Phyllostomidae.
The three known species of vampire bats all seem more similar to one another than to any other species. That suggests that hematophagy evolved only once, and the three species share this common ancestor.
The placement of the three genera of the subfamily Desmodontinae within the New World leaf-nosed bat family Phyllostomidae Gray, 1825, may be summarized as:
subfamily Desmodontinae
genus Desmodus
Desmodus archaeodaptes, extinct,
Desmodus draculae, extinct,
Desmodus rotundus,
Desmodus stocki, extinct.
genus Diphylla
Diphylla ecaudata
genus Diaemus
Diaemus youngi
Evolution
Vampire bats are in a diverse family of bats that consume many food sources, including nectar, pollen, insects, fruit and meat.
The three species of vampire bats are the only mammals that have evolved to feed exclusively on blood (hematophagy) as micropredators, a strategy within parasitism.
Hematophagy is uncommon due to the number of challenges to overcome for success: a large volume of liquid potentially overwhelming the kidneys and bladder, the risk of iron poisoning, and coping with excess protein.
There are multiple hypotheses for how vampire bats evolved.
They evolved from frugivorous bats with sharp teeth specialized for piercing fruit
They initially fed on the ectoparasites of large mammals, and then progressed to feeding on the mammals themselves (similar to red-billed oxpecker feeding behavior)
They initially fed on insects that were attracted to the wounds of animals, and then progressed to feeding on the wounds
They initially preyed on small arboreal vertebrates
They were arboreal omnivores themselves and began ingesting blood and flesh from wound sites of larger animals
They were specialized nectar-feeders that evolved to feed on another type of liquid
The vampire bat lineage diverged from its family 26 million years ago. The hairy-legged vampire bat likely diverged from the other two species of vampire bats 21.7 million years ago.
Because the hairy-legged vampire bat feeds on bird blood and it is the most basal of living vampire bats, it is considered likely that the first vampire bats fed on bird blood as well.
Recent analyses suggest that vampire bats arose from insectivores, which discount the frugivore, carnivore, and nectarivore hypotheses of origin. Within 4 million years of diverging from other Phyllostomidae, vampire bats had evolved all necessary adaptations for blood-feeding, making it one of the fastest examples of natural selection among mammals.
Anatomy and physiology
Unlike fruit bats, the vampire bats have short, conical muzzles. They also lack a nose leaf, instead having naked pads with U-shaped grooves at the tip. The common vampire bat, Desmodus rotundus, also has specialized thermoreceptors on its nose, which aid the animal in locating areas where the blood flows close to the skin of its prey. A nucleus has been found in the brain of vampire bats that has a similar position and similar histology to the infrared receptor of infrared-sensing snakes.
A vampire bat has front teeth that are specialized for cutting and back teeth that are much smaller than in other bats. The inferior colliculus, the part of the bat's brain that processes sound, is well adapted to detecting the regular breathing sounds of sleeping animals that serve as its main food source.
While other bats have almost lost the ability to maneuver on land, vampire bats can walk, jump, and even run by using a unique, bounding gait, in which the forelimbs instead of the hindlimbs are recruited for force production, as the wings are much more powerful than the legs. This ability to run seems to have evolved independently within the bat lineage.
Vampire bats also have a high level of resistance to a group of bloodborne viruses known as endogenous retroviruses, which insert copies of their genetic material into their host's genome.
It was recently discovered that the vampire bat's loss of the REP15 gene allows for enhanced iron secretion in adaptation to the high iron diet.
Vampire bats use infrared radiation to locate blood hotspots on their prey. A recent study has shown that common vampire bats tune a TRP-channel that is already heat-sensitive, TRPV1, by lowering its thermal activation threshold to about . This is achieved through alternative splicing of TRPV1 transcripts to produce a channel with a truncated carboxy-terminal cytoplasmic domain. These splicing events occur exclusively in trigeminal ganglia, and not in dorsal root ganglia, thereby maintaining a role for TRPV1 as a detector of noxious heat in somatic afferents. The only other known vertebrates capable of detecting infrared radiation are boas, pythons and pit vipers, all of which have pit organs.
Ecology and life cycle
Vampire bats tend to live in colonies in almost completely dark places, such as caves, old wells, hollow trees, and buildings. They range in Central to South America and live in arid to humid, tropical and subtropical areas. Vampire bat colony numbers can range from single digits to hundreds in roosting sites. The basic social structure of roosting bats is made of female groups and their offspring, a few adult males, known as "resident males", and a separate group of males, known as "nonresident males". In hairy-legged vampire bats, the hierarchical segregation of nonresident males appears less strict than in common vampire bats. Nonresident males are accepted into the harems when the ambient temperature lowers. This behavior suggests social thermoregulation.
Resident males mate with the females in their harems, and it is less common for outside males to copulate with the females. Female offspring often remain in their natal groups. Several matrilines can be found in a group, as unrelated females regularly join groups. Male offspring tend to live in their natal groups until they are about two years old, sometimes being forcibly expelled by the resident adult males.Vampire bats on average live about nine years when they are in their natural environment in the wild.
Vampire bats form strong bonds with other members of the colony. A related unique adaptation of vampire bats is the sharing of food. A vampire bat can only survive about two days without feeding, yet they cannot be guaranteed of finding food every night. This poses a problem, so when a bat fails to find food, it will often "beg" another bat for food. A "donor" bat may regurgitate a small amount of blood to sustain the other member of the colony. For equally familiar bats, the predictive capacity of reciprocity surpasses that of relatedness. This finding suggests that vampire bats are capable of preferentially aiding their relatives, but that they may benefit more from forming reciprocal, cooperative relationships with relatives and non-relatives alike. Furthermore, donor bats were more likely to approach starving bats and initiate the food sharing. When individuals of a population are lost, bats with a larger number of mutual donors tend to offset their own energetic costs at a higher rate than bats that fed less of the colony before the removal. Individuals that spend their own energy as a social investment of sorts are more likely to thrive, and higher rates of survival incentivize the behavior and reinforce the importance of large social networks in colonies. These findings contradict the harassment hypothesis—which claims that individuals share food in order to limit harassment by begging individuals. All considered, vampire bat research should be interpreted cautiously as much of the evidence is correlational and still requires further testing.
Another ability that some vampire bats possess is identifying and monitoring the positions of conspecifics (individuals of the same species) simply by antiphonal calling. Similar in nature to the sound mother bats make to call to their pups, these calls tend to vary on a bat to bat basis which may help other bats identify individuals both in and outside of their roost.
Vampire bats also engage in social grooming. It usually occurs between females and their offspring, but it is also significant between adult females. Social grooming is mostly associated with food sharing.
Feeding
Vampire bats hunt only when it is fully dark. Like fruit-eating bats, and unlike insectivorous and fish-eating bats, they emit only low-energy sound pulses. The common vampire bat feeds primarily on the blood of mammals (occasionally including humans), whereas both the hairy-legged vampire bat and white-winged vampire bat feed primarily on the blood of birds. Once the common vampire bat locates a host, such as a sleeping mammal, it lands and approaches it on the ground while on all fours. It then likely uses thermoception to identify a warm spot on the skin to bite. They then start to lick the area over and over again to make the place tender so it's easier to bite. Then create a small incision with their teeth and lap up blood from the wound.
Vampire bats, like snakes, have developed highly sensitive thermosensation, with specialized systems for detecting infrared radiation. Snakes co-opt a non-heat-sensitive channel, vertebrate TRPA1 (transient receptor potential cation channel A1), to produce an infrared detector. However, vampire bats tune a channel that is already heat-sensitive, TRPV1, by lowering its thermal activation threshold to about , which allows them to sense the target.
As noted by Arthur M. Greenhall: If there is fur on the skin of the host, the common vampire bat uses its canine and cheek teeth like a barber's blades to shave away the hairs. The bat's razor-sharp upper incisor teeth then make a 7 mm wide and 8 mm deep cut. The upper incisors lack enamel, which keeps them permanently razor sharp. Their teeth are so sharp that even handling their skulls in a museum can result in cuts.
The bat's saliva, left in the victim's resulting bite wound, has a key function in feeding from the wound. The saliva contains several compounds that prolong bleeding, such as anticoagulants that inhibit blood clotting, and compounds that prevent the constriction of blood vessels near the wound.
Digestion
A typical female vampire bat weighs and can consume over 20 grams (1 fluid ounce) of blood in a 20-minute feed. This feeding behavior is facilitated by its anatomy and physiology for rapid processing and digestion of the blood to enable the animal to take flight soon after the feeding.
The stomach and intestine rapidly absorb the water in the blood meal, which is quickly transported to the kidneys, and on to the bladder for excretion. A common vampire bat begins to expel urine within two minutes of feeding.
While shedding much of the blood's liquid facilitates flight takeoff, the bat still has added almost 20–30% of its body weight in blood. To take off from the ground, the bat generates extra lift by crouching and flinging itself into the air. Typically, within two hours of setting out in search of food, the common vampire bat returns to its roost and settles down to spend the rest of the night digesting its meal. Digestion is aided by their microbiome, and their genome protects them against pathogens in the blood. Its stool is roughly the same as that from bats eating fruits or insects.
Metabolism
In a 2024 study published in Biology Letters, researchers explored how vampire bats generate energy from their blood meals, hypothesizing that they metabolize amino acids due to their low-carbohydrate and low-fat diet. The team captured two dozen vampire bats in Belize and fed them cow blood enriched with glycine and leucine. After consumption, the bats were placed on a treadmill for up to 90 minutes, during which breath samples were collected to measure oxygen intake and carbon dioxide output. The results revealed that up to 60% of the bats’ energy production during exercise came from the rapid breakdown of these amino acids, revealing their ability to convert proteins into usable energy within ten minutes. Michael Hiller, a researcher at the LOEWE Center for Translational Biodiversity Genomics in Frankfurt, noted that this rapid metabolization of amino acids is "unparalleled in mammals" and described it as a compelling example of convergent evolution, where both vampire bats and blood-feeding insects developed similar strategies to adapt to their extreme diets.
This metabolic specialization presents drawbacks, as vampire bats have diminished their ability to store alternative energy sources, rendering them susceptible to starvation if they experience prolonged periods without feeding. To counteract this vulnerability, vampire bats engage in reciprocal altruism, regurgitating blood to assist conspecifics in need.
Human health
Rabies
Rabies can be transmitted to humans and other animals by vampire bat bites. Since dogs are now widely immunized against rabies, the number of human rabies transmissions by vampire bats exceeds those by dogs in Latin America, with 55 documented cases in 2005. The risk of infection to the human population is less than to livestock exposed to bat bites. Various estimates of the prevalence of rabies in bat populations have been made; it has been estimated that less than 1% of wild bats in regions where rabies is endemic are infected with the virus at any given time. Bats that are infected may be clumsy, disoriented, and unable to fly.
Anticoagulant drug
The unique properties of vampire bat saliva have found some positive use in medicine.
Various studies published in Stroke: Journal of the American Heart Association on a genetically engineered drug called desmoteplase which uses the anticoagulant properties of the saliva of Desmodus rotundus found that it increased blood flow in stroke patients.
| Biology and health sciences | Bats | null |
430347 | https://en.wikipedia.org/wiki/Peony | Peony | The peony or paeony () is any flowering plant in the genus Paeonia, the only genus in the family Paeoniaceae. Peonies are native to Asia, Europe, and Western North America. Scientists differ on the number of species that can be distinguished, ranging from 25 to 40, although the current consensus describes 33 known species. The relationships between the species need to be further clarified.
Most are herbaceous perennial plants tall, but some are woody shrubs tall. They have compound, deeply lobed leaves and large, often fragrant flowers, in colors ranging from purple and pink to red, white or yellow, in late spring and early summer. The flowers have a short blooming season, usually lasting for only 7–10 days.
Peonies are popular garden plants in temperate regions. Herbaceous peonies are also sold as cut flowers on a large scale, although they generally are only available in late spring and early summer.
Description
Morphology
All Paeoniaceae are herbaceous perennials or deciduous shrubs, with thick storage roots and thin roots for gathering water and minerals. Some species are caespitose (tufted), because the crown produces adventitious buds, while others have stolons. They have rather large compound leaves without glands and stipules, and with anomocytic stomata. In the woody species the new growth emerges from scaly buds on the previous flush or from the crown of the rootstock. The large bisexual flowers are mostly single at the end of the stem. In P. emodi, P. lactiflora, P. veitchii and many of the cultivars these contributed to, few additional flowers develop in the axils of the leaves. Flowers close at night or when the sky is overcast. Each flower is subtended by a number of bracts, that may form a sort of involucre, has 3-7 tough free sepals and mostly 5–8, but occasionally up to 13 free petals. These categories however are intergrading, making it difficult to assign some of them, and the number of these parts may vary. Within are numerous (50–160) free stamens, with anthers fixed at their base to the filaments, and are sagittate in shape, open with longitudinal slits at the outer side and free pollen grains which have three slits or pores and consist of two cells. Within the circle of stamens is a more or less prominent, lobed disc, which is presumed not to excrete nectar. Within the disk is a varying number (1-15) of separate carpels, which have a very short style and a decurrent stigma. Each of these develops into a dry fruit (which is called a follicle), which opens with a lengthwise suture and each of which contains one or a few large fleshy seeds. The annual growth is predetermined: if the growing tip of a shoot is removed, no new buds will develop that season.
Phytochemistry
Over 262 compounds have been obtained so far from the plants of Paeoniaceae. These include monoterpenoid glucosides, flavonoids, tannins, stilbenoids, triterpenoids, steroids, paeonols, and phenols. In vitro biological activities include antioxidant, antitumor, antipathogenic, immunomodulative, cardiovascular-system-protective activities and central-nervous-system activities.
Paeoniaceae are dependent on C3 carbon fixation. They contain ellagic acid, myricetin, ethereal oils and flavones, as well as crystals of calcium oxalate. The wax tubules that are formed primarily consist of palmitone (the ketone of palmitic acid).
Genome
The basic chromosome number is five. About half of the species of the section Paeonia however is tetraploid (4n=20), particularly many of those in the Mediterranean region. Both allotetraploids and autotetraploids are known, and some diploid species are also of hybrid origin.
Taxonomy
The family name "Paeoniaceae" was first used by Friedrich K.L. Rudolphi in 1830, following a suggestion by Friedrich Gottlieb Bartling that same year. The family had been given other names a few years earlier. The composition of the family has varied, but it has always consisted of Paeonia and one or more genera that are now placed in Ranunculales. It has been widely believed that Paeonia is closest to Glaucidium, and this idea has been followed in some recent works. Molecular phylogenetic studies, however, have demonstrated conclusively that Glaucidium belongs in the family Ranunculaceae, order Ranunculales, but that Paeonia belongs in the unrelated order Saxifragales. The genus Paeonia consists of about 35 species, assigned to three sections: Moutan, Onaepia and Paeoniae. The section Onaepia only includes P. brownii and P. californica. The section Moutan is divided into P. delavayi and P. ludlowii, together making up the subsection Delavayanae, and P. cathayana, P. decomposita, P. jishanensis, P. osti, P. qiui and P. rockii which constitute the subsection Vaginatae. P. suffruticosa is a cultivated hybrid swarm, not a naturally occurring species.
The remainder of the species belongs to the section Paeonia, which is characterised by a complicated reticulate evolution. Only about half of the (sub)species is diploid, the other half tetraploid, while some species both have diploid and tetraploid populations. In addition to the tetraploids, are some diploid species also likely the result of hybridisation, or nothospecies. Known diploid taxa in the Paeonia-section are P. anomala, P. lactiflora, P. veitchii, P. tenuifolia, P. emodi, P. broteri, P. cambedessedesii, P. clusii, P. rhodia, P. daurica subsps. coriifolia, daurica, macrophylla and mlokosewitschii. Tetraploid taxa are P. arietina, P. officinalis, P. parnassica, P. banatica, P. russi, P. peregrina, P. coriacea, P. mascula subsps. hellenica and mascula, and P. daurica subsps. tomentosa and wittmanniana. Species that have both diploid and tetraploid populations include P. clusii, P. mairei and P. obovata. P. anomala was proven to be a hybrid of P. lactiflora and P. veitchii, although being a diploid with 10 chromosomes. P. emodi and P. sterniana are diploid hybrids of P. lactiflora and P. veitchii too, and radically different in appearance. P. russi is the tetraploid hybrid of diploid P. lactiflora and P. mairei, while P. cambedessedesii is the diploid hybrid of P. lactiflora, likely P. mairei, but possibly also P. obovata. P. peregrina is the tetraploid hybrid of P. anomala and either P. arietina, P. humilis, P. officinalis, P. parnassica or less likely P. tenuifolia, or one of their (now extinct) common ancestors. P. banatica is the tetraploid hybrid of P. mairei and one of this same group. P. broteri, P. coriacea, P. clusii, P. rhodia, P. daurica subsp. mlokosewitschi, P. mascula subsp. hellenica and ssp. mascula, and P. daurica subsp. wittmanniana are all descendants of hybrids of P. lactiflora and P. obovata.
Phylogeny
Recent genetic analyses relate the monogeneric family Paeoniaceae to a group of families with woody species in the order Saxifragales. This results in the following relationship tree. One dissertation suggests the section Onaepia branches off earliest, but a later publication of the same author and others suggests the Moutan-section splits off first. Within that section P. ludlowii and P. delavayi are more related to each other than to any other species.
Species
Herbaceous species (about 30 species)
Paeonia algeriensis (Algerian peony)
Paeonia anomala
Paeonia turcica (Turkish peony)
Paeonia arietina (ram's horn peony)
Paeonia broteri (Brotero's peony)
Paeonia brownii (Brown's peony, native peony, or western peony)
Paeonia californica (California peony or wild peony)
Paeonia cambessedesii (Majorcan peony or Balearic peony)
Paeonia clusii
subsp. clusii
subsp. rhodia
Paeonia coriacea (Andalusian peony)
Paeonia corsica (Corsican peony)
Paeonia daurica (Crimean peony)
subsp. coriifolia
subsp. daurica
subsp. macrophylla
subsp. mlokosewitschii (Caucasian peony or golden peony)
subsp. tomentosa
subsp. velebitensis
subsp. wittmanniana (Wittmann's peony)
Paeonia emodi (Himalayan peony)
Paeonia intermedia (Intermediate peony or Altay peony)
Paeonia kesrouanensis (Keserwan peony)
Paeonia lactiflora (Chinese peony, Chinese herbaceous peony, or common garden peony)
Paeonia mairei
Paeonia mascula (Balkan peony, wild peony, or male peony)
subsp. mascula
subsp. bodurii
subsp. hellenica
subsp. russoi
Paeonia obovata (woodland peony)
subsp. willmottiae
Paeonia officinalis (European peony, common peony, or garden peony; type species)
Paeonia parnassica (Greek peony)
Paeonia peregrina (Kosovar peony or Romanian peony)
Paeonia sterniana
Paeonia tenuifolia (steppe peony or fern leaf peony)
Paeonia veitchii (Veitch's peony)
Woody species, tree peonies (about 8 species)
Paeonia decomposita
Paeonia delavayi (Delavay's tree peony or Dian peony)
Paeonia jishanensis (Jishan peony)
Paeonia ludlowii (Tibetan tree peony or Ludlow's tree peony)
Paeonia ostii (Osti's peony)
Paeonia qiui (Qiu's peony)
Paeonia rockii (Rock's peony; synonym Paeonia suffruticosa subsp. rockii (Chinese tree peony, known as "moutan (moutan peony)" in China))
Distribution
The genus Paeonia naturally occurs in the temperate and cold areas of the Northern Hemisphere. The section Moutan, which includes all woody species, is restricted in the wild to Central and Southern China, including Tibet. The section Onaepia consist of two herbaceous species and is present in the West of North-America, P. brownii between southern British Columbia and the Sierra Nevada in California and eastward to Wyoming and Utah, while P. californica is limited to the coastal mountains of Southern and Central California.
The section Paeonia, which comprises all other herbaceous species, occurs in a band stretching roughly from Morocco and Spain to Japan. One species of the section Paeonia, P. anomala, has by far the largest distribution, which is also north of the distribution of the other species: from the Kola peninsula in North-West Russia, to Lake Baikal in Siberia and South to the Tien Shan Mountains of Kazakhstan. The rest of the section concentrates around the Mediterranean, and in Asia.
The species around the Mediterranean include Paeonia algeriensis that is an endemic of the coastal mountains of Algeria, P. coriacea in the Rif Mountains and Andalusia, P. cambessedesii on Majorca, P. russoi on Corsica, Sardinia and Sicily, P. corsica on Corsica, Sardinia, the Ionian islands and in western Greece, P. clusii subsp. clusii on Crete and Karpathos, and subsp. rhodia on Rhodes, P. kesrouanensis in the Western Taurus Mountains, P. arietina from the Middle Taurus Mountains, P. broteri in Andalucia, P. humilis from Andalucia to the Provence, P. officinalis from the South of France, through Switzerland to the Middle of Italy, P. banatica in western Romania, northern Serbia and Slovenia and in southern Hungary, P. peregrina in Albania, western Bulgaria, northern Greece, western Romania, Serbia, Montenegro and Bosnia, while P. mascula has a large distribution from Catalonia and southern France to Israel and Turkey.
Between the two concentrations, the subspecies of Paeonia daurica occur, with subspecies velebitensis in Croatia, and daurica in the Balkans and Crimea, while the other subspecies coriifolia, macrophylla, mlokosewitschii, tomentosa and wittmanniana are known from the Caucasus, Kaçkar and Alborz Mountains.
Paeonia emodi occurs in the western Himalayas between Pakistan and western Nepal, P. sterniana is an endemic of southeastern Tibet, P. veitchii grows in Central China (Qinghai, Ningxia, Gansu, Shaanxi, Shanxi, Sichuan and the eastern rim of Tibet), like P. mairei (Gansu, Guizhou, Hubei, Shaanxi, Sichuan, and Yunnan), while P. obovata grows in warm-temperate to cold China, including Manchuria, Korea, Japan, Far Eastern Russia (Primorsky Krai) and on Sakhalin, and P. lactiflora occurs in Northern China, including Manchuria, Japan, Korea, Mongolia, Russia (Far East and Siberia).
Distributional history
The species of the section Paeonia have a disjunct distribution, with most of the species occurring in the Mediterranean, while many others occur in eastern Asia. Genetic analysis has shown that all Mediterranean species are either diploid or tetraploid hybrids that resulted from the crossbreeding of species currently limited to eastern Asia. The large distance between the ranges of the parent species and the nothospecies suggest that hybridisation already occurred relatively long ago. It is likely that the parent species occurred in the same region when the hybrids arose, and were later exterminated by successive Pleistocene glaciations, while the nothospecies remained in refugia to the South of Europe. During their retreat, P. lactiflora and P. mairei likely became sympatric and so produced the Himalayan nothospecies P. emodi and P. sterniana.
Cultivation
Ancient Chinese texts mention the peony was used for flavoring food. Peonies have been used and cultivated in China since early history. Ornamental cultivars were created from plants cultivated for medicine in China as of the sixth and seventh century. Peonies became particularly popular during the Tang dynasty, when they were grown in the imperial gardens. In the tenth century the cultivation of peonies spread through China, and the seat of the Song dynasty, Luoyang, was the centre for its cultivation, a position it still holds today.
A second centre for peony cultivation developed during the Qing dynasty in Cáozhōu, now known as He Ze. Both cities still host annual peony exhibitions and state-funded peony research facilities. Before the tenth century, P. lactiflora was introduced in Japan, and over time many varieties were developed both by self fertilisation and crossbreeding, particularly during the eighteenth to twentieth centuries (middle Edo to early Shōwa periods). During the 1940s Toichi Itoh succeeded in crossing tree peonies and herbaceous peonies and so created a new class of so-called intersectional hybrids. Although P. officinalis and its cultivars were grown in Europe from the fifteenth century on, originally also for medicinal purposes, intensive breeding started only in the nineteenth century when P. lactiflora was introduced from its native China to Europe. The tree peony was introduced in Europe and planted in Kew Gardens in 1789. The main centre of peony breeding in Europe has been in the United Kingdom, and particularly France. Here, breeders like Victor Lemoine and François Félix Crousse selected many new varieties, mainly with P. lactiflora, such as "Avant Garde" and "Le Printemps". The Netherlands is the largest peony cut flower producing country with about 50 million stems each year, with "Sarah Bernhardt" dominating the sales with over 20 million stems. An emerging source of peonies in mid to late summer is the Alaskan market. Unique growing conditions due to long hours of sunlight create availability from Alaska when other sources have completed harvest.
Plant growth habits
While the peony takes several years to re-establish itself when moved, it blooms annually for decades once it has done so.
Peonies tend to attract ants to the flower buds. This is due to the nectar that forms on the outside of the flower buds, and is not required for the plants' own pollination or other growth. The presence of ants is thought to provide some deterrence to other harmful insects though, so the production of ant-attracting nectar is plausibly a functional adaptation. Ants do not harm the plants.
Peony species come in two distinct growth habits, while hybrid cultivars in addition may occupy an intermediate habit.
herbaceous: During summer, renewal buds develop on the underground stem (the "crown"), particularly at the foot of the current season's annual shoots. These renewal buds come in various sizes. Large buds will grow into stems the following growing season, but smaller buds remain dormant. The primordia for the leaves can already be found in June, but the flower only starts differentiating in October, as the annual shoots die down, completing its development in December, when sepals, petals, stamens and pistils are all recognisable.
tree: During the summer, large buds develop at the tip of the annual growth and near its foot. In the autumn, the leaves are shed, and the new stems become woody and are perennial.
Itoh (or "Intersectional"): In 1948 horticulturist Toichi Itoh from Tokyo used pollen from the yellow tree peony "Alice Harding" to fertilize the herbaceous P. lactiflora "Katoden", which resulted in a new category of peonies, the Itoh or intersectional cultivars. These are herbaceous, have leaves like tree peonies, with many large flowers from late spring to early autumn, and good peony wilt resistance. Some of the early Itoh cultivars are "Yellow Crown", "Yellow Dream", "Yellow Emperor" and "Yellow Heaven".
Flower types
Seven types of flower are generally distinguished in cultivars of herbaceous peonies.
single: a single or double row of broad petals encircle fertile stamens, carpels visible.
Japanese: a single or double row of broad petals encircle somewhat broadened staminodes, may carry pollen along the edges, carpels visible.
anemone: a single or double row of broad petals encircle narrow incurved petal-like staminodes; fertile stamens are absent, carpels visible.
triple: the flower consists of triple row of broad petals that broaden and overlap each other.
semi-double: a single or double row of broad petals encircles further broad petals intermingled with stamens.
bomb: a single row of broad petals encircles a shorter dense pompon of narrower petals.
double: the flower consists of many broad petals only, including those which likely are altered stamens and carpels.
Propagation
Herbaceous and Itoh peonies are propagated by root division, and sometimes by seed. Tree peonies can be propagated by grafting, division, seed, and from cuttings, although root grafting is most common commercially.
Herbaceous peonies such as Paeonia lactiflora, will die back to ground level each autumn. Their stems will reappear the following spring. However tree peonies, such as Paeonia suffruticosa, are shrubbier. They produce permanent woody stems that will lose their leaves in winter but the stem itself remains intact above ground level.
Cultivars
The numerous peony hybrids and cultivars have gained the Royal Horticultural Society's Award of Garden Merit, including:
'Bartzella', a double yellow-flowered Itoh (intersectional) peony
'Coral Charm', a semi-double salmon-pink-flowered herbaceous peony
Paeonia × festiva 'Rubra Plena', a bomb red-flowered herbaceous peony
Paeonia × lemoinei 'High Noon', a semi-double yellow-flowered tree peony
The American Peony Society is the International Cultivar Registration Authority for the genus, and accepts over 7,000 registered cultivars.
Uses
The herb known as Paeonia, in particular the root of P. lactiflora (Bai Shao, Radix Paeoniae Lactiflorae), has been used frequently in traditional medicines of Korea, China and Japan. In Japan, Paeonia lactiflora used to be called ebisugusuri ("foreign medicine"). Pronunciation of 牡丹 (peony) in Japan is "botan." In kampo, the Japanese adaptation of Chinese medicine, its root was used as a treatment for convulsions. It is also cultivated as a garden plant. In Japan Paeonia suffruticosa is called the "King of Flowers" and Paeonia lactiflora is called the "Prime Minister of Flowers."
In China, the fallen petals of Paeonia lactiflora are parboiled and sweetened as a tea-time delicacy. Peony water, an infusion of peony petals, was used for drinking in the Middle Ages. The petals may be added to salads or to punches and lemonades.
Culture
The peony is among the longest-used flowers in Eastern culture. Along with the plum blossom, it is a traditional floral symbol of China, where the Paeonia suffruticosa is called 牡丹 (mǔdān). It is also known as 富貴花 (fùguìhuā) "flower of riches and honour" or 花王 (huawang) "king of the flowers", and is used symbolically in Chinese art.
In 1903, the Qing dynasty declared the peony as the national flower. Currently, the Republic of China government in Taiwan designates the plum blossom as the national flower, while the People's Republic of China government has no legally designated national flower. In 1994, the peony was proposed as the national flower after a nationwide poll, but the National People's Congress failed to ratify the selection. In 2003, another selection process was initiated, but no choice has been made to date.
The ancient Chinese city Luoyang has a reputation as a cultivation centre for the peonies. Throughout Chinese history, peonies in Luoyang have been said to be the finest in the country. Dozens of peony exhibitions and shows are still held there annually.
The Greek doctor Dioscorides named aglaophotis, an herb supposedly capable of warding off certain evils, as a member of the peony family.
In the Middle Ages, peonies were often painted with their ripe seed-capsules, since it was the seeds, not the flowers, which were medically significant. Ancient superstition dictated that great care be taken not to be seen by a woodpecker while picking the plant's fruit, or the bird might peck out one's eyes.
The red flowers of the species Paeonia peregrina are important in Serbian folklore. Known as Kosovo peonies (, ), they are said to represent the blood of Serbian warriors who died in the Battle of Kosovo.
In 1957, the Indiana General Assembly passed a law to make the peony the state flower of Indiana, a title which it holds to this day. It replaced the zinnia, which had been the state flower since 1931.
Mischievous nymphs were said to hide in the petals of the peony, giving it the meaning of Shame or Bashfulness in the Language of Flowers.
Peonies are a common subject in tattoos, often used along with koi-fish. The popular use of peonies in Japanese tattoo was inspired by the ukiyo-e artist Utagawa Kuniyoshi's illustrations of Suikoden, a classical Chinese novel. His paintings of warrior-heroes covered in pictorial tattoos included lions, tigers, dragons, koi fish, and peonies, among other symbols. The peony became a masculine motif, associated with a devil-may-care attitude and disregard for consequence.
Famous painters of peonies have included Conrad Gessner (ca. 1550) and Auguste Renoir in 1879. Paeonia officinalis can be found in the altar picture of Maria im Rosenhag by Schongauer in the former Dominican Church in Colmar. The Italian Jesuit, painter and architect Giuseppe Castiglione (1688-1766), who worked at the court of the Qianlong Emperor in the Qing dynasty, also painted peonies.
| Biology and health sciences | Saxifragales | Plants |
430368 | https://en.wikipedia.org/wiki/Swan%20goose | Swan goose | The swan goose (Anser cygnoides) is a large goose with a natural breeding range in inland Mongolia, Northeast China, and the Russian Far East. It is migratory and winters mainly in central and eastern China. Vagrant birds are encountered in Japan and Korea (where it used to winter in numbers when it was more common), and more rarely in Kazakhstan, Laos, coastal Siberia, Taiwan, Thailand and Uzbekistan.
While uncommon in the wild, this species has been extensively domesticated, when it is known as Chinese goose. Introduced and feral populations of its domestic breeds occur in many places outside its natural range. The wild form is also kept in collections, and escapes are not unusual amongst feral flocks of other Anser and Branta geese.
Taxonomy
The swan goose was formally described in 1758 by the Swedish naturalist Carl Linnaeus in the tenth edition of his Systema Naturae under the binomial name Anas cygnoid. The original spelling of "cygnoid." (with a fullstop, suggesting an abbreviation) is widely treated as an orthographic error to be corrected to "cygnoides", but some authors regard cygnoid (without the fullstop) to be the correct original spelling which should be used, citing the comparable case, universally accepted, of the spelling of the name Estrilda astrild. As of early 2025, this dispute over the spelling remains unresolved.
It is now one of 11 species placed in the genus Anser that was described by the French zoologist Mathurin Jacques Brisson in 1760. The specific epithet combines the Latin cygnus meaning "swan" with Ancient Greek -oidēs meaning "resembling". The species is monotypic, with no subspecies are recognised.
Description
The swan goose is large and long-necked for its genus, wild birds being long (the longest Anser goose) and weighing or more (the second-heaviest Anser, after the greylag goose A. anser). The sexes are similar, although the male is larger, with a proportionally longer bill and neck; in fact the largest females are barely as large as the smallest males. Typical measurements of the wing are in males, in females; the bill is about long in males and in females. The tarsus of males measures around . The wingspan of adult geese is .
The upperparts are greyish-brown, with thin light fringes to the larger feathers and a maroon hindneck and cap (reaching just below the eye). The remiges are blackish, as are the entire underwing and the white-tipped rectrices, while the upper- and undertail coverts are white. A thin white stripe surrounds the bill base. Apart from darker streaks on the belly and flanks, the underside is pale buff, being especially light on the lower head and foreneck which are sharply delimited against the maroon. In flight, the wings appear dark, with no conspicuous pattern. Uniquely among its genus, the long, heavy bill is completely black; the legs and feet, on the other hand, are orange as in most of its relatives. The eyes' irides are maroon. Juveniles are duller than adult birds, and lack the white bill base and dark streaks on the underside.
The voice is a loud, drawn-out and ascending honking aang. As a warning call, a similar but more barking honk is given two or three times in short succession.
The karyotype of the swan goose is 2n=80, consisting of four pairs of macrochromosomes, 35 pairs of microchromosomes, and a pair of sex chromosomes. The two largest macrochromosome pairs as well as the Z (female) chromosome are submetacentric, while the third-largest chromosome pair is acrocentric and the fourth-largest is metacentric. The W chromosomes are acrocentric too, as are the larger microchromosomes, the smaller ones probably being telocentric. Compared to the greylag goose, there seems to have been some rearrangement on the fourth-largest chromosome pair.
Ecology
It inhabits steppe to taiga and mountain valleys near freshwater, grazing on plants such as sedges (Cyperaceae), grasses (Poaceae) and water plants, and rarely swimming. It forms small flocks outside the breeding season. In the winter, it grazes on plains and stubble fields, sometimes far from water. Birds return from the winter quarters around April, and the breeding season starts soon thereafter. It breeds as single pairs or loose groups near marshes and other wetlands, with nesting activity starting about May. The clutch is usually 5–6 but sometimes up to 8 eggs, which are laid in a shallow nest made from plants, placed directly on the ground, often on a small knoll to keep it dry. The precocial young hatch after about 28 days and become sexually mature at 2–3 years of age. Around late August/early September, the birds leave for winter quarters, where they gather in small groups to moult their worn plumage.
The swan goose was uplisted from Near Threatened to Vulnerable on the IUCN Red List in 1992 and further to Endangered in 2000, as its population is declining due to habitat loss and excessive hunting and (particularly on the Sanjiang Plain in China) egg collecting. But new research has shown it to be not as rare as it was believed, and consequently, it was downlisted to Vulnerable status again in 2008. Still, less than 500 pairs might remain in Russia, while in Mongolia numbers are unknown, though about 1,000 were seen at Ögii Lake in 1977. Important wintering locations in China are Lake Dongting, Lake Poyang, the Yancheng Coastal Wetlands and other locations around the lower Yangtze River, where some 60,000 individuals may be found each year – though this may be almost the entire world population. Until the 1950s, the species wintered in small numbers (up to about 100 birds annually) in Japan, but habitat destruction has driven them away.
Domestication
Though the majority of domestic geese are descended from the greylag goose (A. anser), two breeds are direct descendants of the swan goose, the Chinese goose and the African goose. These breeds have been domesticated since at least the mid-18th century – perhaps even (in China) since around 1000 BC. They vary considerably from their wild parent in appearance, temperament, and ability to produce meat and eggs; the most conspicuous feature is the prominent bill knob and upright posture.
Charles Darwin studied goose breeds as part of his work on the theory of evolution. He noted that the external differences between Chinese geese and breeds descended from the Greylag goose belied a rather close relationship:
"The hybrids from the common and Chinese geese (A. cygnoides), species which are so different that they are generally ranked in distinct genera, have often bred in this country with either pure parent, and in one single instance they have bred inter se."
Conservation
The species is currently classified as an endangered species by the IUCN based on ongoing population declines and range losses, exacerbated by recent poor breeding success and unsustainable levels of hunting. Total population was estimated as 36–43,500 individuals in 2023.
Gallery
| Biology and health sciences | Anseriformes | Animals |
430380 | https://en.wikipedia.org/wiki/Orthopedic%20surgery | Orthopedic surgery | Orthopedic surgery or orthopedics (alternative spelling orthopaedics) is the branch of surgery concerned with conditions involving the musculoskeletal system. Orthopedic surgeons use both surgical and nonsurgical means to treat musculoskeletal trauma, spine diseases, sports injuries, degenerative diseases, infections, tumors, and congenital disorders.
Etymology
Nicholas Andry coined the word in French as , derived from the Ancient Greek words ("correct", "straight") and ("child"), and published Orthopedie (translated as Orthopædia: Or the Art of Correcting and Preventing Deformities in Children) in 1741. The word was assimilated into English as orthopædics; the ligature æ was common in that era for ae in Greek- and Latin-based words. As the name implies, the discipline was initially developed with attention to children, but the correction of spinal and bone deformities in all stages of life eventually became the cornerstone of orthopedic practice.
Differences in spelling
As with many words derived with the "æ" ligature, simplification to either "ae" or just "e" is common, especially in North America. In the US, the majority of college, university, and residency programmes, and even the American Academy of Orthopaedic Surgeons, still use the spelling with the digraph ae, though hospitals usually use the shortened form. Elsewhere, usage is not uniform; in Canada, both spellings are acceptable; "orthopaedics" is the normal spelling in the UK in line with other fields which retain "ae".
History
Early orthopedics
Many developments in orthopedic surgery have resulted from experiences during wartime. On the battlefields of the Middle Ages, the injured were treated with bandages soaked in horses' blood, which dried to form a stiff, if unsanitary, splint.
Originally, the term orthopedics meant the correcting of musculoskeletal deformities in children. Nicolas Andry, a professor of medicine at the University of Paris, coined the term in the first textbook written on the subject in 1741. He advocated the use of exercise, manipulation, and splinting to treat deformities in children. His book was directed towards parents, and while some topics would be familiar to orthopedists today, it also included 'excessive sweating of the palms' and freckles.
Jean-André Venel established the first orthopedic institute in 1780, which was the first hospital dedicated to the treatment of children's skeletal deformities. He developed the club-foot shoe for children born with foot deformities and various methods to treat curvature of the spine.
Advances made in surgical technique during the 18th century, such as John Hunter's research on tendon healing and Percival Pott's work on spinal deformity steadily increased the range of new methods available for effective treatment. Robert Chessher, a pioneering British orthopedist, invented the double-inclined plane, used to treat lower-body bone fractures, in 1790. Antonius Mathijsen, a Dutch military surgeon, invented the plaster of Paris cast in 1851. Until the 1890s, though, orthopedics was still a study limited to the correction of deformity in children. One of the first surgical procedures developed was percutaneous tenotomy. This involved cutting a tendon, originally the Achilles tendon, to help treat deformities alongside bracing and exercises. In the late 1800s and first decades of the 1900s, significant controversy arose about whether orthopedics should include surgical procedures at all.
Modern orthopedics
Examples of people who aided the development of modern orthopedic surgery were Hugh Owen Thomas, a surgeon from Wales, and his nephew, Robert Jones. Thomas became interested in orthopedics and bone-setting at a young age, and after establishing his own practice, went on to expand the field into the general treatment of fracture and other musculoskeletal problems. He advocated enforced rest as the best remedy for fractures and tuberculosis, and created the so-called "Thomas splint" to stabilize a fractured femur and prevent infection. He is also responsible for numerous other medical innovations that all carry his name: Thomas's collar to treat tuberculosis of the cervical spine, Thomas's maneuvere, an orthopedic investigation for fracture of the hip joint, the Thomas test, a method of detecting hip deformity by having the patient lying flat in bed, and Thomas's wrench for reducing fractures, as well as an osteoclast to break and reset bones.
Thomas's work was not fully appreciated in his own lifetime. Only during the First World War did his techniques come to be used for injured soldiers on the battlefield. His nephew, Sir Robert Jones, had already made great advances in orthopedics in his position as surgeon-superintendent for the construction of the Manchester Ship Canal in 1888. He was responsible for the injured among the 20,000 workers, and he organized the first comprehensive accident service in the world, dividing the 36-mile site into three sections, and establishing a hospital and a string of first-aid posts in each section. He had the medical personnel trained in fracture management. He personally managed 3,000 cases and performed 300 operations in his own hospital. This position enabled him to learn new techniques and improve the standard of fracture management. Physicians from around the world came to Jones' clinic to learn his techniques. Along with Alfred Tubby, Jones founded the British Orthopedic Society in 1894.
During the First World War, Jones served as a Territorial Army surgeon. He observed that treatment of fractures both, at the front and in hospitals at home, was inadequate, and his efforts led to the introduction of military orthopedic hospitals. He was appointed Inspector of Military Orthopedics, with responsibility for 30,000 beds. The hospital in Ducane Road, Hammersmith, became the model for both British and American military orthopedic hospitals. His advocacy of the use of Thomas splint for the initial treatment of femoral fractures reduced mortality of open fractures of the femur from 87% to less than 8% in the period from 1916 to 1918.
The use of intramedullary rods to treat fractures of the femur and tibia was pioneered by Gerhard Küntscher of Germany. This made a noticeable difference to the speed of recovery of injured German soldiers during World War II and led to more widespread adoption of intramedullary fixation of fractures in the rest of the world. Traction was the standard method of treating thigh bone fractures until the late 1970s, though, when the Harborview Medical Center group in Seattle popularized intramedullary fixation without opening up the fracture.
The modern total hip replacement was pioneered by Sir John Charnley, expert in tribology at Wrightington Hospital, in England in the 1960s. He found that joint surfaces could be replaced by implants cemented to the bone. His design consisted of a stainless steel, one-piece femoral stem and head, and a polyethylene acetabular component, both of which were fixed to the bone using PMMA (acrylic) bone cement. For over two decades, the Charnley low-friction arthroplasty and its derivative designs were the most-used systems in the world. This formed the basis for all modern hip implants.
The Exeter hip replacement system (with a slightly different stem geometry) was developed at the same time. Since Charnley, improvements have been continuous in the design and technique of joint replacement (arthroplasty) with many contributors, including W. H. Harris, the son of R. I. Harris, whose team at Harvard pioneered uncemented arthroplasty techniques with the bone bonding directly to the implant.
Knee replacements, using similar technology, were started by McIntosh in rheumatoid arthritis patients and later by Gunston and Marmor for osteoarthritis in the 1970s, developed by John Insall in New York using a fixed bearing system, and by Frederick Buechel and Michael Pappas using a mobile bearing system.
External fixation of fractures was refined by American surgeons during the Vietnam War, but a major contribution was made by Gavril Abramovich Ilizarov in the USSR. He was sent, without much orthopedic training, to look after injured Russian soldiers in Siberia in the 1950s. With no equipment, he was confronted with crippling conditions of unhealed, infected, and misaligned fractures. With the help of the local bicycle shop, he devised ring external fixators tensioned like the spokes of a bicycle. With this equipment, he achieved healing, realignment, and lengthening to a degree unheard of elsewhere. His Ilizarov apparatus is still used today as one of the distraction osteogenesis methods.
Modern orthopedic surgery and musculoskeletal research have sought to make surgery less invasive and to make implanted components better and more durable. On the other hand, since the emergence of the opioid epidemic, orthopedic surgeons have been identified as one of the highest prescribers of opioid medications. Decreasing prescription of opioids while still providing adequate pain control is a development in orthopedic surgery.
Training
In the United States, orthopedic surgeons have typically completed four years of undergraduate education and four years of medical school and earned either a Doctor of Medicine (MD) or Doctor of Osteopathic Medicine (DO) degree. Subsequently, these medical school graduates undergo residency training in orthopedic surgery. The five-year residency is a categorical orthopedic surgery training.
Selection for residency training in orthopedic surgery is very competitive. Roughly 700 physicians complete orthopedic residency training per year in the United States. About 10% of current orthopedic surgery residents are women; about 20% are members of minority groups. Around 20,400 actively practicing orthopedic surgeons and residents are in the United States. According to the latest Occupational Outlook Handbook (2011–2012) published by the United States Department of Labor, 3–4% of all practicing physicians are orthopedic surgeons.
Many orthopedic surgeons elect to do further training, or fellowships, after completing their residency training. Fellowship training in an orthopedic sub-specialty is typically one year in duration (sometimes two) and sometimes has a research component involved with the clinical and operative training. Examples of orthopedic subspecialty training in the United States are:
Foot and ankle surgery
Hand and upper extremities
Hip and knee surgery
Orthopedic oncologist
Orthopedic trauma
Osseointegration
Pediatric orthopedics
Shoulder and elbow
Spine surgery
Surgical sports medicine
Total joint reconstruction (arthroplasty)
These specialized areas of medicine are not exclusive to orthopedic surgery. For example, hand surgery is practiced by some plastic surgeons, and spine surgery is practiced by most neurosurgeons. Additionally, foot and ankle surgery is also practiced by doctors of podiatric medicine (DPM) in the United States. Some family practice physicians practice sports medicine, but their scope of practice is nonoperative.
After completion of specialty residency or registrar training, an orthopedic surgeon is then eligible for board certification by the American Board of Medical Specialties or the American Osteopathic Association Bureau of Osteopathic Specialists. Certification by the American Board of Orthopedic Surgery or the American Osteopathic Board of Orthopedic Surgery means that the orthopedic surgeon has met the specified educational, evaluation, and examination requirements of the board. The process requires successful completion of a standardized written examination followed by an oral examination focused on the surgeon's clinical and surgical performance over a 6-month period. In Canada, the certifying organization is the Royal College of Physicians and Surgeons of Canada; in Australia and New Zealand, it is the Royal Australasian College of Surgeons.
In the United States, specialists in hand surgery and orthopedic sports medicine may obtain a certificate of added qualifications in addition to their board primary certification by successfully completing a separate standardized examination. No additional certification process exists for the other subspecialties.
Practice
According to applications for board certification from 1999 to 2003, the top 25 most common procedures (in order) performed by orthopedic surgeons are:
Knee arthroscopy and meniscectomy
Shoulder arthroscopy and decompression
Carpal tunnel release
Knee arthroscopy and chondroplasty
Removal of support implant
Knee arthroscopy and anterior cruciate ligament reconstruction
Knee replacement
Repair of femoral neck fracture
Repair of trochanteric fracture
Debridement of skin/muscle/bone/ fracture
Knee arthroscopy repair of both menisci
Hip replacement
Shoulder arthroscopy/distal clavicle excision
Repair of rotator cuff tendon
Repair fracture of radius/ulna
Laminectomy
Repair of ankle fracture (bimalleolar type)
Shoulder arthroscopy and debridement
Lumbar spinal fusion
Repair fracture of the distal part of radius
Low back intervertebral disc surgery
Incise finger tendon sheath
Repair of ankle fracture (fibula)
Repair of femoral shaft fracture
Repair of trochanteric fracture
A typical schedule for a practicing orthopedic surgeon involves 50–55 hours of work per week divided among clinic, surgery, various administrative duties, and possibly teaching and/or research if in an academic setting. According to the American Association of Medical Colleges in 2021, the average work week of an orthopedic surgeon was 57 hours. This is a very low estimation however, as research derived from a 2013 survey of orthopedic surgeons who self identified as "highly successful" due to their prominent positions in the field indicated average work weeks of 70 hours or more.
Arthroscopy
The use of arthroscopic techniques has been particularly important for injured patients. Arthroscopy was pioneered in the early 1950s by Masaki Watanabe of Japan to perform minimally invasive cartilage surgery and reconstructions of torn ligaments. Arthroscopy allows patients to recover from the surgery in a matter of days, rather than the weeks to months required by conventional, "open" surgery; it is a very popular technique. Knee arthroscopy is one of the most common operations performed by orthopedic surgeons today, and is often combined with meniscectomy or chondroplasty. The majority of upper-extremity outpatient orthopedic procedures are now performed arthroscopically.
Arthroplasty
Arthroplasty is an orthopedic surgery where the articular surface of a musculoskeletal joint is replaced, remodeled, or realigned by osteotomy or some other procedure. It is an elective procedure that is done to relieve pain and restore function to the joint after damage by arthritis (rheumasurgery) or some other type of trauma. As well as the standard total knee replacement surgery, the unicompartmental knee replacement, in which only one weight-bearing surface of an arthritic knee is replaced, may be performed, but it bears a significant risk of revision surgery. Joint replacements are used for other joints, most commonly the hip or shoulder.
A post-surgical concern with joint replacements is wear of the bearing surfaces of components. This can lead to damage to the surrounding bone and contribute to eventual failure of the implant. The plastic chosen is usually ultra-high-molecular-weight polyethylene, which can also be altered in ways that may improve wear characteristics. The risk of revision surgery has also been shown to be associated with surgeon volume.
Epidemiology
Between 2001 and 2016, the prevalence of musculoskeletal procedures drastically increased in the U.S., from 17.9% to 24.2% of all operating-room (OR) procedures performed during hospital stays.
In a study of hospitalizations in the United States in 2012, spine and joint procedures were common among all age groups except infants. Spinal fusion was one of the five most common OR procedures performed in every age group except infants younger than 1 year and adults 85 years and older. Laminectomy was common among adults aged 18–84 years. Knee arthroplasty and hip replacement were in the top five OR procedures for adults aged 45 years and older.
| Biology and health sciences | Fields of medicine | Health |
430383 | https://en.wikipedia.org/wiki/Urtica%20dioica | Urtica dioica | Urtica dioica, often known as common nettle, burn nettle, stinging nettle (although not all plants of this species sting) or nettle leaf, or just a nettle or stinger, is a herbaceous perennial flowering plant in the family Urticaceae. Originally native to Europe, much of temperate Asia and western North Africa, it is now found worldwide. The species is divided into six subspecies, five of which have many hollow stinging hairs called trichomes on the leaves and stems, which act like hypodermic needles, injecting histamine and other chemicals that produce a stinging sensation upon contact ("contact urticaria", a form of contact dermatitis).
The plant has a long history of use as a source for traditional medicine, food, tea, and textile raw material in ancient (such as Saxon) and modern societies.
Description
Urtica dioica is a dioecious, herbaceous, and perennial plant. It grows to tall in the summer and dying down to the ground in winter. It has widely spreading rhizomes and stolons, which are bright yellow, as are the roots. The soft, green leaves are long and are borne oppositely on an erect, wiry, green stem. The leaves have a strongly serrated margin, a cordate base, and an acuminate tip with a terminal leaf tooth longer than adjacent laterals. It bears small, greenish or brownish, numerous flowers in dense axillary inflorescences.
The leaves and stems are very hairy with non-stinging hairs, and in most subspecies, also bear many stinging hairs (trichomes or spicules), whose tips come off when touched, transforming the hair into a needle that can inject several chemicals causing a painful sting or paresthesia, giving the species its common names: stinging nettle, burn-nettle, burn-weed, or burn-hazel.
Taxonomy
Credit for the scientific naming of Urtica dioica is given to Carl Linnaeus, who published it in Species Plantarum in 1753. The taxonomy of Urtica species is confused, and sources are likely to use a variety of systematic names for these plants. Until 2014 there was broad consensus that the nettles native to the Americas, now classified as Urtica gracilis, were subspecies of U. dioica. However, in that year the paper "Weeding the Nettles II" was published in the journal Phytotaxa demonstrating the genetic distinctness of New World nettles. As of 2023 Plants of the World Online (POWO) recognizes U. gracilis as a distinct species while the USDA Natural Resources Conservation Service PLANTS database continues to list it as U. dioica subsp. gracilis, as does the Flora of North America.
As of 2023 POWO recognizes 11 subspecies or varieties of U. dioica:
Urtica dioica subsp. afghanica Chrtek, from southwestern and central Asia, sometimes has stinging hairs or is sometimes hairless.
Urtica dioica subsp. dioica (European stinging nettle), from Europe, Asia, and northern Africa, has stinging hairs.
Urtica dioica subsp. gansuensis C.J.Chen, from eastern Asia (China), has stinging hairs.
Urtica dioica var. glabrata (Clem.) Asch. & Graebn.
Urtica dioica var. hispida (Lam. ex DC.) Tausch ex Ott
Urtica dioica var. holosericea Fr.
Urtica dioica subsp. kurdistanica Chrtek
Urtica dioica subsp. pubescens(Ledeb.) Domin, in many sources as U. dioica subsp. galeopsifolia (fen nettle or stingless nettle), from Europe, does not have stinging hairs.
Urtica dioica var. sarmatica Zapał.
Urtica dioica subsp. sondenii(Simmons) Hyl.
Urtica dioica subsp. subinermis (R.Uechtr.) Weigend
Etymology
Urtica is derived from a Latin word meaning 'sting'.
Dioica (δίοικος) is derived from Greek, meaning 'of two houses' (having separate staminate and pistillate plants; dioecious).
Distribution and habitat
U. dioica is considered to be native to Europe, much of temperate Asia and western North Africa. It is abundant in northern Europe and much of Asia, usually found in the countryside. It is less widespread in southern Europe and north Africa, where it is restricted by its need for moist soil, but is still common. It has been introduced to many other parts of the world. In North America, it is widely distributed in Canada and the United States, where it is found in every province and state except for Hawaii, and also can be found in northernmost Mexico. It grows in abundance in the Pacific Northwest, especially in places where annual rainfall is high. The European subspecies has been introduced into Australia, North America and South America.
In Europe, nettles have a strong association with human habitation and buildings. The presence of nettles may indicate the site of a long-abandoned building, and can also indicate soil fertility. Human and animal waste may be responsible for elevated levels of phosphate and nitrogen in the soil, providing an ideal environment for nettles.
Ecology
Nettles are the larval food plant for several species of butterflies, such as the peacock butterfly, comma (Polygonia c-album), and the small tortoiseshell. It is also eaten by the larvae of some moths including angle shades, buff ermine, dot moth, the flame, the gothic, grey chi, grey pug, lesser broad-bordered yellow underwing, mouse moth, setaceous Hebrew character, and small angle shades. The roots are sometimes eaten by the larva of the ghost moth (Hepialus humuli).
It is a known host to the pathogenic fungus Phoma herbarum.
Stinging nettle is particularly found as an understory plant in wetter environments, but it is also found in meadows. Although nutritious, it is not widely eaten by either wildlife or livestock, presumably because of the sting. It spreads by abundant seeds and also by rhizomes, and is often able to survive and re-establish quickly after fire.
Cultivation
Field
Sowing and planting
Three cultivation techniques can be used for the stinging nettle: 1) direct sowing, 2) growing seedlings in nurseries with subsequent transplantation and 3) vegetative propagation via stolons or head cuttings.
Direct sowing: The seedbed should have a loose and fine structure, but should be reconsolidated using a packer roller imminently prior to sowing. Sowing time can be either in autumn or in spring. Seed density should be 6 kilograms/hectare with row spacing of and 42–50 cm in autumn and spring, respectively. The disadvantage of direct sowing is that it usually leads to incomplete plant coverage. This drawback can be mitigated by covering the seedbed with a transparent perforated foil in order to improve seed germination. Further, weed control can be problematic as the stinging nettle has a slow seedling development time.
Growing seedlings: For this technique pre-germinated seeds are sown between mid-/end-February and beginning of April and grown in nurseries. Seedlings are grown in tuffs with 3–5 plants/tuff and a seed density of 1.2–1.6 kg/1000 tuffs. Faster germination is achieved by alternating high temperature during daytime (30 °C for 8 h) and lower temperature during nighttime (20 °C for 16 h). Before transplanting, the seedlings should be fertilized and acclimated to cold temperatures. Transplantation should start around Mid-April with row spacing of and plant spacing within rows of 25–30 cm.
Vegetative propagation: Stolons (with several buds) of 10 cm should be planted from mid-April in a depth of . Head cuttings are grown in nurseries starting between mid-May and mid-June. Growing tips with two leaf pairs are cut from the mother plant and treated with root-growth inducing hormones. Transplantation can be delayed in comparison to the growing seedling technique.
Greenhouse
The stinging nettle can also be grown in controlled-environment agriculture systems, such as soil-less medium cultivations or aeroponics, which may achieve higher yields, standardize quality, and reduce harvesting costs and contamination.
Sting and treatment
Urtica dioica produces its inflammatory effect on skin (a stinging, burning sensation often called "contact urticaria") both by impaling the skin via spicules – causing mechanical irritation – and by biochemical irritants, such as histamine, serotonin, and acetylcholine, among other chemicals. Anti-itch drugs, usually in the form of creams containing antihistamines or hydrocortisone, may provide relief from nettle dermatitis. The term contact urticaria has a wider use in dermatology, involving dermatitis caused by various skin irritants and pathogens.
Docks, especially Rumex obtusifolius (the broad-leaf dock) often grow in similar environments to stinging nettles and are regarded as a folk remedy to counteract the sting of a nettle, although there is no evidence of any chemical effect. It may be that the act of rubbing a dock leaf against a nettle sting acts as a distracting counterstimulation, or that belief in the dock's effect provides a placebo effect.
Uses
Culinary
U. dioica has a flavour similar to spinach when cooked. Young plants are harvested by many Native American communities and are cooked and eaten in spring when other food plants are scarce. Soaking stinging nettles in water or cooking removes the stinging chemicals from the plant, which allows them to be handled and eaten without injury. After the stinging nettle enters its flowering and seed-setting stages, the leaves develop gritty particles called cystoliths. Many sources claim consumption of these can irritate the kidneys and urinary tract, but there is no medical evidence to support this claim. Cystoliths are made of calcium carbonate, and will not dissolve when boiled. Leaves harvested post-flowering must have their cystoliths broken down by acid, as in the fermentation process. In its peak season, nettle contains up to 25% protein, dry weight, which is high for a leafy green vegetable. The leaves are also dried and may then be used to make a herbal tea, as can also be done with the nettle's flowers.
Nettles can be used in a variety of recipes, such as polenta, pesto, and purée. Nettle soup is a common use of the plant, particularly in Northern and Eastern Europe.
Nettles are sometimes used in cheesemaking, such as for Cornish Yarg and as a flavouring in varieties of Gouda.
Nettles are used in Albania, Montenegro, Serbia, North Macedonia and Bosnia and Herzegovina as part of the dough filling for the börek pastry. The top baby leaves are selected and simmered, and then mixed with other ingredients such as herbs and rice, before being used as a filling between dough layers. Similarly, in Greece the tender leaves are often used, after simmering, as a filling for hortopita, which is similar to spanakopita, but with wild greens rather than spinach for filling.
Young nettles can also be used to make an alcoholic drink.
Competitive eating
In the United Kingdom, an annual World Nettle Eating Championship draws thousands of people to Dorset, where competitors attempt to eat as much of the raw plant as possible. Competitors are given stalks of the plant, from which they strip the leaves and eat them. Whoever strips and eats the most stinging nettle leaves in a fixed time is the winner. The competition dates back to 1986, when two neighbouring farmers attempted to settle a dispute about which had the worst infestation of nettles, and one of them said to the other, "I'll eat any nettle of yours that's longer than mine."
Traditional medicine
As Old English , nettle is one of the nine plants invoked in the pagan Anglo-Saxon Nine Herbs Charm, recorded in 10th-century traditional medicine. Nettle was believed to be a galactagogue – a substance that promotes lactation. Urtication, or flogging with nettles, is the process of deliberately applying stinging nettles to the skin to provoke inflammation. An agent thus used was considered to be a rubefacient (something that causes redness), used as a folk remedy for treating rheumatism. A study undertaken in 2000 showed that nettles were an effective therapy in relieving the pain of arthritis.
Chastisement
In indigenous justice systems in Ecuador, urtication was used as punishment for severe crimes in 2010. The sentenced perpetrator of a crime was flogged with stinging nettle, in public, naked, whilst being showered with freezing cold ice water.
Textiles and fibre
Nettle stems contain a bast fibre that has been traditionally used for the same purposes as linen and is produced by a similar retting process. Unlike cotton, nettles grow easily without pesticides. The fibres are coarser, however.
Historically, nettles have been used to make clothing for almost 3,000 years, as ancient nettle textiles from the Bronze Age have been found in Denmark. It is widely believed that German Army uniforms were almost all made from nettle during World War I due to a shortage of cotton, although there is little evidence to support this. More recently, companies in Austria, Germany, and Italy have started to produce commercial nettle textiles.
The fibre content in nettle shows a high variability and reaches from below 1% to 17%. Under middle-European conditions, stems yield typically between 45 and 55 dt / ha (decitons per hectare), which is comparable to flax stem yield. Due to the variable fibre content, the fibre yields vary between 0.2 and 7 dt / ha, but the yields are normally in the range between 2 and 4 dt / ha. Fibre varieties are normally cloning varieties and therefore planted from vegetative propagated plantlets. Direct seeding is possible, but leads to great heterogeneity in maturity.
Nettles may be used as a dye-stuff, producing yellow from the roots, or yellowish green from the leaves.
Feed
Nutrient contents
Fresh leaves contain approximately 82.4% water, 17.6% dry matter, 5.5% protein, 0.7 to 3.3% fat, and 7.1% carbohydrates. Mature leaves contain about 40% α- linolenic acid, a valuable omega-3 acid. For exact fatty acid contents see Table 1. Seeds contain much more fatty acid than leaves.
Minerals (Ca, K, Mg, P, Si, S, Cl) and trace elements (Ti, 80 ppm, Mn, Cu, Fe) contents depend mostly on the soil and the season.
Carotenoids can be found primarily in the leaves, where different forms of lutein, xanthophyll and carotene are present (Table 2). Some carotenes are precursors of vitamin A (retinol), their retinol equivalents RE or retinol activity equivalents per g dry weight are 1.33 for mature leaves and 0.9 for young leaves. Nettle contains much less carotenes and retinol than carrots, which contain 8.35 RE per g fresh weight. Depending on the batch and the leaf and stem content, nettle contains only traces of zeaxanthin or between 20–60 mg/kg of dry matter. Nettle contains ascorbic acid (vitamin C), riboflavin (vitamin B2), pantothenic acid, vitamin K1 and tocopherols (vitamin E). The highest vitamin contents can be found in the leaves.
Poultry: Egg yolk colouring in laying hens
In laying hens, nettle can be used as an egg yolk colourant instead of artificial pigments or other natural pigments (derived from marigold for yellow). Nettle has high carotenoid contents, especially lutein, and zeaxanthin, of which lutein and zeaxanthin act as yellow pigments. Feeding as little as 6.25 g dry nettle per kg feed is as effective as the synthetic pigments to colour the egg yolk. Feeding nettle has no detrimental effect on the performance of the laying hens or on the general quality of eggs.
Ruminants
Ruminants avoid fresh stinging nettles; however, if the nettles are wilted or dry, voluntary intake can be high.
Use in agriculture / horticulture
In the European Union, nettle extract can be used as an insecticide, fungicide, and acaricide under Basic Substance regulations. As an insecticide, nettle extract can be used for the control of codling moth, diamondback moth, and spider mites. As a fungicide, it can be used for the control of Pythium root rot, powdery mildew, early blight, late blight, Septoria blight, Alternaria leaf spot, and grey mould.
Gardening
Nettles have a number of uses in the vegetable garden, including the potential for encouraging beneficial insects. Since nettles prefer to grow in phosphorus-rich and nitrogen rich soils that have recently been disturbed (and thus aerated), the growth of nettles is an indicator that an area has high fertility (especially phosphate and nitrate), and thus is an indicator to gardeners as to the quality of the soil.
Nettles contain nitrogenous compounds, so are used as a compost activator or can be used to make a liquid fertilizer, which although low in phosphate, is useful in supplying magnesium, sulphur, and iron. They are also one of the few plants that can tolerate, and flourish in, soils rich in poultry droppings.
The stinging nettle is the red admiral caterpillar's primary host plant and can attract migrating red admiral butterflies to a garden. U. dioica can be a troubling weed, and mowing can increase plant density. Regular and persistent tilling will greatly reduce its numbers, and the use of herbicides such as and glyphosate are effective control measures.
In culture
In Great Britain and Ireland, U. dioica and the annual nettle Urtica urens are the only common stinging plants and have found a place in several figures of speech in the English language. Shakespeare's Hotspur urges that "out of this nettle, danger, we pluck this flower, safety" (Henry IV, Part 1, Act II Scene 3). The figure of speech "to grasp the nettle" probably originated from Aesop's fable "The Boy and the Nettle". In Seán O'Casey's Juno and the Paycock, one of the characters quotes Aesop "Gently touch a nettle and it'll sting you for your pains/Grasp it as a lad of mettle and soft as silk remains". The metaphor may refer to the fact that if a nettle plant is grasped firmly rather than brushed against, it does not sting so readily, because the hairs are crushed down flat and do not penetrate the skin so easily.
In the German language, the idiom , or to sit in nettles, means to get into trouble. In Germanic mythology, the God of thunder, Thor, was associated with nettles, and that's where the saying "lightning won't strike into nettles" comes from. The idiom is used in Croatian, Hungarian, Serbian, and many other Indo-European languages. In Dutch, a means a predicament. In French, the idiom (do not push granny into the nettles) means that we should be careful not to abuse a situation. The name urticaria for hives comes from the Latin name of nettle (, from urere, to burn).
The English word 'nettled', meaning irritated or angry, is derived from 'nettle'.
There is a common idea in Great Britain that the nettle was introduced by the Romans, but Plant Atlas 2020 treats it as native. The idea of its introduction was mentioned by William Camden in his book Britannia of 1586. However, in 2011, an early Bronze Age burial cist on Whitehorse Hill, Dartmoor, Devon, was excavated. The cist dated from between 1730 and 1600 BC. It contained various high value beads as well as fragments of a sash made from nettle fibre. It is possible that the sash was traded from mainland Europe, but perhaps more probable that it was locally made.
| Biology and health sciences | Rosales | Plants |
430783 | https://en.wikipedia.org/wiki/Athens%20Metro | Athens Metro | The Athens Metro () is a rapid transit system serving the Athens urban area in Greece. Line 1 opened as a single-track conventional steam railway in 1869 and was electrified in 1904. Beginning in 1991, Elliniko Metro S.A. constructed and extended Lines 2 and 3.
It has significantly changed Athens by providing a much-needed solution to the city's traffic and air pollution problem, as well as revitalising many of the areas it serves.
Extensions of existing lines are under development or tender, like the Line 2 extension to Ilion where tender started in 2023, as well as a new Line 4, whose central section began construction in October 2021.
The Athens Metro is actively connected with the other means of public transport, such as buses, trolleys, the Athens Tram and the Athens Suburban Railway. The Athens Metro is hailed for its modernity (mainly the newer lines 2, 3), and many of its stations feature works of art, exhibitions and displays of the archaeological remains found during its construction. Photography and video-taking is permitted across the whole network and street photographers often work in Athens Metro. This was the only metro system in Greece, before the Thessaloniki Metro began operations on 30th November 2024.
History
Piraeus–Kifissia Railway (Line 1)
Until 28 January 2000, Line 1 was the only rapid-transit line in Athens. The Athens and Piraeus Railway Company (SAP) opened a steam single-track mixed cargo and passenger railway line on 27 February 1869 and was run between and . It was electrified in 1904. On 4 February 1885 Lavrion Square-Strofyli steam narrow gauge single-track mixed cargo and passenger railway line opened and was run at the time from Attiki Square to Kifissia through Iraklio. These originally mixed cargo and passenger railway lines gradually merged and converted to a rapid-transit system. The section between Kifissia and Strofyli was abandoned.
From 1869 to 1926 the line was operated by SAP. From 1926 to 1976 the line was operated by Hellenic Electric Railways (EIS). In 1976 the EIS was nationalized and renamed Athens-Piraeus Electric Railway Company (ISAP), which continued to operate what became line 1 of the Athens Metro until 16 June 2011.
1990s projects
Since the current Line 1 opened, the government has proposed many expansions to the subway network, including a 1963 plan for a fourteen-line subway network. Construction of Lines 2 and 3 began in November 1992 to decrease traffic congestion and improve Athens' air quality by reducing its smog level. Both lines were constructed underground. Lines 2 and 3, built by Attiko Metro S.A. and operated until 2011 by Attiko Metro Operations Company, are known respectively as the red and blue lines and were inaugurated in January 2000. Line 3 was extended to the Eleftherios Venizelos International Airport in summer 2004, and Line 2 was extended to Anthoupoli and Elliniko in 2013.
Consolidation
Until 17 June 2011, the operational management of the Athens Metro network was similar to that of the London Underground network before the creation of the London Passenger Transport Board and the absorption of the Metropolitan Railway on 1 July 1933. The Greek government attempted to absorb ISAP into Attiko Metro operation company under Law 2669/1998 so the latter would be responsible for the whole network, but this initiative failed. Athens Metro operations were consolidated when the Greek government enacted Law 3920/2011, replacing AMEL, ISAP and Tram S.A. with Urban Rail Transport S.A. (STASY S.A.) (), a subsidiary of OASA S.A. (Athens Urban Transport Organisation S.A.).
Timeline
Infrastructure
Lines and stations
The Athens Metro consists of three lines totalling and 66 stations: Line 1 (Green) is long with 24 stations, Line 2 (Red) is long with 20 stations, and Line 3 (Blue) is long with 24 stations. STASY owns and operates 62 of the 66 stations: three other stations (, and ) belong to GAIAOSE and the station belongs to the operator of the Athens International Airport.
The system has five interchanges, at , , , and , allowing all three to interchange with each other at least once. Each line also has at least one connection with the Athens Suburban Railway, and the Athens Tram.
Line 2 is entirely underground. Line 1 is mostly overground, with an underground section spanning between the Monastiraki and Attiki stations, and an additional underground station (Kato Patisia) in central Athens. Line 3 is mostly underground; Trains that run an overground route are only those with the airport as final destination. The overground section of Line 3, east of the tunnel portal near , is open. In the tunnel sections up and down lines share a common tunnel, except for approaches to stations with an island platform (such as Egaleo). Train maintenance facilities are located at Attiki, Faliro, Irini, Piraeus, Kifissia and Thissio for Line 1, and Doukissis Plakentias, Eleonas and Sepolia for Lines 2 and 3.
The Athens Metro's three lines carried approximately 1,353,000 passengers daily in 2010.
A network map of the Athens Metro system, that includes the three current lines, the under construction line 4, the tramway, the suburban railway and all the future under design extensions.
Rolling stock
The network uses standard gauge electric trains which in most places run on 750 V DC third rail, but the section of Line 3 running to the airport requires trains which can use overhead lines of 25 kV AC, 50 Hz.
The Athens Metro classifies rolling stock by "batch" for Line 1 and "generation" for Lines 2 and 3 because ISAP and AMEL used different classification systems for rolling stock before consolidation. Six types of rolling stock operate on the network, all equipped with third rail current collection systems; however, only seven second-generation trains have the necessary overhead line equipment to serve Line 3 from to .
The eighth batch (introduced in 1983) is the oldest rolling stock in passenger service, while the third generation (introduced in 2013) is the latest rolling stock in passenger service. The eighth- and tenth-batch stock is externally similar, but the former has split-flap headsigns in Johnston typeface and a cream-and-green interior colour scheme. An extensive refurbishment programme is in progress for the 8th batch (as of 2023), and to cover for trains undergoing refurbishment, up to five 1st generation Line 2/3 trains have been borrowed to operate on Line 1.
Line 1 halfsets have driving cabs at both ends, unlike the Line 2/3 halfsets which have a driving cab at the outer ends, but only basic driving apparatus for shunting purposes only at the inner ends; thus, they can only operate on their own inside depots.
Line 1
Lines 2 & 3
Signalling
Line 1 uses two-aspect red/green home signals, yellow/green distant signals and a passenger information system (PIS). The current system replaced 1950s-era semaphore signals. The automatic train protection (ATP) system of Line 1 was fully installed in 2023 which replaced the previous Indusi system.
Lines 2 and 3 use the Alstom automatic train supervision system (ATS) and a passenger information system (PIS). Two-aspect red/white colour signals are used at points and junctions only.
Fares
Fares are prepaid, either as short term tickets valid for 90 minutes, 24 hours, three days, five days, or as long term tickets. As of September 2020, there are two types of fare products, the ATH.ENA Ticket and ATH.ENA Card, both of which are validated using a contactless system (by scanning the ticket or card at the electronic validating machines). The tickets are valid on all modes of public transport in Athens except on trains and buses to the airport. Passengers cannot buy a fare on board the bus. To travel to or from the airport, passengers may buy a one-way ticket for €9 or a 3-day ticket for €20 which also includes unlimited local trips and a return trip to the airport. Arrival at the airport without having paid the appropriate fare will incur a €72 fine, reduced to €36 if paid within 10 days. Term tickets are available in 30, 90, 180, and 365 day periods and are available only with a personalized ATH.ENA Card. Reduced fares are available for university students, seniors, disabled and persons under 18. During a fare control the passengers that are entitled to a reduced fare have to show ID card, student card or passport. Children under the age of 6 are entitled to travel for free with all means of transportation. On buses and trams the ticket or card must be validated only when entering the vehicle/car by scanning the ticket at the electronic validating machines. At metro or Suburban Railway stations, the ticket or card must be validated at the electronic gates when entering and exiting the station.
Archaeological excavations and exhibits
During construction of the metro tunnels, artifacts of archaeological interest were discovered and rescue archaeology was employed. Teams of archaeologists worked ahead of, then with, engineers for six years, protecting and recording archaeological finds (streets, houses, cemeteries, sanctuaries, public workshops, foundry pits, kilns, aqueducts, wells, cisterns, drains and sewage tunnels). This afforded new insight into the city's ancient topography, through unprecedented infrastructure development combined with the study and preservation of archaeological data. Exhibitions of ancient artifacts or replicas are found at a number of metro stations, including Monastiraki, Akropoli and Syntagma.
Future
The Athens Metro masterplan, as presented in October 2022, consists of the following projects:
*The current Kifissia terminal will be demolished and rebuilt as an underground station.
**The Development Plan refers it as Line 4 branch but there are unofficial plans that this branch is part of the future Line 5.
If and when these projects are completed, the Athens Metro is expected to reach in length and serve a total of 110 stations by 2040.
Line 4
A fourth line is planned for the Athens Metro and it has been incorporated in the roadmap for the development and expansion of public transport in Athens since 2005. The new line in its totality will extend over a length of , adding thirty five (35) new stations to the Athens Metro system. The cost of the entire project is estimated at 3.3 billion EUR. The recommendation is for lighter rolling stock than the type used in existing lines of Athens Metro which would operate automatically without a driver. In November 2020, Alstom was chosen to supply the line with 20 4-car automated Metropolis trains, operated under Urbalis 400 signalling system.
The first phase of Line 4 will be between Alsos Veikou and Goudi stations, predicting fifteen (15) new stations and a length of of new track. An invitation to tender for the construction of the first phase of Line 4 was issued in September 2018. The construction started in mid to late-2021 and is scheduled to be completed in 2029 or 2030. The estimated cost for constructing the first phase of the new line is 1.51 billion EUR. Currently, the project of the first phase is considered to follow a PPP scheme which might be extended for constructing the whole new line. An alternative solution is a mixed funding between the EIB and the Greek State. It is also a high-profile candidate project to be included in the Juncker Plan of EU that will include also the second phase of Line 4 of Athens Metro.
The European Investment Bank is allocating a €730 million loan over 30 years to finance the building of the first segment of the Line 4 metro, which will connect Alsos Veikou and Goudi. The initiative will also assist Athens by reducing the number of private automobiles on the road by 53 000, resulting in 318 tonnes fewer CO2 released daily.
Long-term plan
On 15 November 2008, Greek newspaper Ta Nea reported that the Greek government was considering a circular line from Ano Ilisia to Faros, via and , as part of a "" network. This proposal evolved to form part of what is now the long-term Athens Metro Future Regulatory Plan (or the Souflias plan) on 13 April 2009, which called for an network of eight lines and 200 stations.
The Souflias plan was last revised in January 2012, and saw limited activity until October 2020, when Elliniko Metro announced that they were reconsidering some extensions from the plan, including the extension of Line 1 from to Nea Erythraia, the extensions of Line 2 to and Glyfada, Line 6 from Melissia to Perama, and Line 7 from Haidari to Kalamaki. In December 2021, a part of the southern branch of Line 6 was reconsidered as a branch of Line 1 from to the SNFCC in Kallithea, with intermediate stations at Hamosternas, Plateia Davaki, and Lofos Filaretou.
| Technology | Europe_2 | null |
430790 | https://en.wikipedia.org/wiki/Gauge%20boson | Gauge boson | In particle physics, a gauge boson is a bosonic elementary particle that acts as the force carrier for elementary fermions. Elementary particles whose interactions are described by a gauge theory interact with each other by the exchange of gauge bosons, usually as virtual particles.
Photons, W and Z bosons, and gluons are gauge bosons. All known gauge bosons have a spin of 1 and therefore are vector bosons. For comparison, the Higgs boson has spin zero and the hypothetical graviton has a spin of 2.
Gauge bosons are different from the other kinds of bosons: first, fundamental scalar bosons (the Higgs boson); second, mesons, which are composite bosons, made of quarks; third, larger composite, non-force-carrying bosons, such as certain atoms.
Gauge bosons in the Standard Model
The Standard Model of particle physics recognizes four kinds of gauge bosons: photons, which carry the electromagnetic interaction; W and Z bosons, which carry the weak interaction; and gluons, which carry the strong interaction.
Isolated gluons do not occur because they are colour-charged and subject to colour confinement.
Multiplicity of gauge bosons
In a quantized gauge theory, gauge bosons are quanta of the gauge fields. Consequently, there are as many gauge bosons as there are generators of the gauge field. In quantum electrodynamics, the gauge group is U(1); in this simple case, there is only one gauge boson, the photon. In quantum chromodynamics, the more complicated group SU(3) has eight generators, corresponding to the eight gluons. The three W and Z bosons correspond (roughly) to the three generators of SU(2) in electroweak theory.
Massive gauge bosons
Gauge invariance requires that gauge bosons are described mathematically by field equations for massless particles. Otherwise, the mass terms add non-zero additional terms to the Lagrangian under gauge transformations, violating gauge symmetry. Therefore, at a naïve theoretical level, all gauge bosons are required to be massless, and the forces that they describe are required to be long-ranged. The conflict between this idea and experimental evidence that the weak and strong interactions have a very short range requires further theoretical insight.
According to the Standard Model, the W and Z bosons gain mass via the Higgs mechanism. In the Higgs mechanism, the four gauge bosons (of SU(2)×U(1) symmetry) of the unified electroweak interaction couple to a Higgs field. This field undergoes spontaneous symmetry breaking due to the shape of its interaction potential. As a result, the universe is permeated by a non-zero Higgs vacuum expectation value (VEV). This VEV couples to three of the electroweak gauge bosons (W, W and Z), giving them mass; the remaining gauge boson remains massless (the photon). This theory also predicts the existence of a scalar Higgs boson, which has been observed in experiments at the LHC.
Beyond the Standard Model
Grand unification theories
The Georgi–Glashow model predicts additional gauge bosons named X and Y bosons. The hypothetical X and Y bosons mediate interactions between quarks and leptons, hence violating conservation of baryon number and causing proton decay. Such bosons would be even more massive than W and Z bosons due to symmetry breaking. Analysis of data collected from such sources as the Super-Kamiokande neutrino detector has yielded no evidence of X and Y bosons.
Gravitons
The fourth fundamental interaction, gravity, may also be carried by a boson, called the graviton. In the absence of experimental evidence and a mathematically coherent theory of quantum gravity, it is unknown whether this would be a gauge boson or not. The role of gauge invariance in general relativity is played by a similar symmetry: diffeomorphism invariance.
W′ and Z′ bosons
W′ and Z′ bosons refer to hypothetical new gauge bosons (named in analogy with the Standard Model W and Z bosons).
| Physical sciences | Bosons | null |
430984 | https://en.wikipedia.org/wiki/Silkie | Silkie | The Silkie (also known as the Silky or Chinese silk chicken) is a Chinese breed of chicken named for its atypically fluffy plumage, which is said to feel like silk and satin. The breed has several other unusual qualities, such as black skin and bones, blue earlobes, and five toes on each foot, whereas most chickens have only four. They are often exhibited in poultry shows, and also appear in various colors. In addition to their distinctive physical characteristics, Silkies are well known for their calm and friendly temperament. It is among the most docile of poultry. Hens are also exceptionally broody, and care for young well. Although they are fair layers themselves, laying only about three eggs a week, they are commonly used to hatch eggs from other breeds and bird species due to their broody nature. Silkie chickens have been bred to have a wide variety of colors which include but are not limited to: Black, Blue, Buff, Partridge, Splash, White, Lavender, Paint and Porcelain.
History
It is unknown exactly where or when these fowl with their singular combination of attributes first appeared, but the most well documented point of origin is ancient China. Other places in Southeast Asia have been named as possibilities, such as India and Java. The earliest surviving Western written account of Silkies comes from Marco Polo, who wrote of a "furry" chicken in the 13th century during his travels in Asia. In 1598, Ulisse Aldrovandi, a writer and naturalist at the University of Bologna, Italy, published a comprehensive treatise on chickens which is still read and admired today. In it, he mentions "wool-bearing chickens" and ones "clothed with hair like that of a black cat".
Silkies most likely made their way to the West via the Silk Route and maritime trade. The breed was recognized officially in North America with acceptance into the Standard of Perfection in 1874. Once Silkies became more common in the West, many myths were perpetuated about them. Early Dutch breeders told buyers they were the offspring of chickens and rabbits, while sideshows promoted them as having actual mammalian fur.
In the 21st century, Silkies are one of the most popular and ubiquitous ornamental breeds of chicken. They are often kept as ornamental fowl or pet chickens by backyard keepers, and are also commonly used to incubate and raise the offspring of other chickens and waterfowl like ducks, geese and game birds such as quail and pheasants.
Characteristics
Silkies are considered a bantam breed in some countries, but this varies according to region and many breed standards class them officially as large fowl; the bantam Silkie is actually a separate variety most of the time. Almost all North American strains of the breed are bantam-sized, but in Europe the standard-sized is the original version. However, even standard Silkies are relatively small chickens, with the males weighing only , and females weighing . The American Standard of Perfection calls for males that are , and females that are .
Silkie plumage was once unique among chicken breeds, however in recent years silkie feathering has been developed in several breeds, mostly notably the Chabo, where it is now standardised in Britain and the Netherlands. It has been compared to silk, and to fur. The overall result is a soft, fluffy appearance. Their feathers lack functioning barbicels, and are thus similar to down on other birds. This characteristic leaves Silkies unable to fly.
Silkies appear in two distinct varieties: bearded and non-bearded. Bearded Silkies have an extra muff of feathers under the beak area that covers the earlobes. They also are separated according to color. Colors of Silkie recognized for competitive showing include black, blue, splash, lavender, buff, grey, partridge, and white. Alternative hues, such as cuckoo, mottled, chocolate, mauve, mille fleur, and red, are in various stages of development and/or awaiting official recognition. The standards of perfection call for all Silkies to have a small walnut-shaped comb, dark wattles, and turquoise-blue earlobes. In addition to these defining characteristics, Silkies have five toes on each foot. Other breeds which exhibit this rare trait include the Dorking, Faverolles, Houdan, and Sultan.
All Silkies have black or bluish skin, bones and grayish-black meat; they are in the group of Chinese fowls known by the Chinese language name of wu gu ji (烏骨雞), meaning 'black-boned chicken'. More specifically, the Silkie breed itself is named Taihe wu ji (泰和乌鸡), 'black-boned chicken from Taihe'. Other wu gu ji may not share characteristics of the Taihe breed, such as the mulberry comb, white fur, blue ears, and polydactyly.
Melanism which extends beyond the skin into an animal's connective tissue is a rare trait, and in chickens it is caused by fibromelanosis, which is a rare mutation believed to have begun in Asia. The Silkie and several other breeds descended from Asian stock possess the mutation. Disregarding color, the breed does not generally produce as much as the more common meat breeds of chicken.
Bantams
In the American Standard of Perfection, the standard male weight for the bantam Silkie is and for the female, . The Australian Poultry Standard and British Poultry Standard call for Silkie bantams much smaller; in the Australian, the standard weights are 680 g (25 oz) for males and 570 g (20 oz) for females. The British standard weight for bantam Silkies is 600 g (22 oz) for males, and 500 g for females (18 oz).
Polydactyly
Silkies are also known for their polydactyly, usually manifesting as an additional 1–2 digits in the foot. The genetic cause of this extra digit formation has been shown to be a SNP in a regulator of the SHH gene, called the ZPA Regulatory Sequence (ZRS). This causes ectopic SHH expression in the anterior of the developing limb bud, leading to increased tissue growth and digits. While the feet of the Silkie display polydactyly, the wings have the standard tridactyly (three digit) arrangement. The Japanese Silkie initially develops additional digits in the wing as an embryo, but these are lost prior to hatching. The genetic cause behind Silkie polydactyly differs from those that cause polydactyly in the Dorking chicken breed, which is due to ectopic FGF4 expression in the AER, with ectopic SHH a secondary effect.
Use
Silkies lay a fair number of eggs, ranging from white to cream or light tan, but production is often interrupted due to their extreme tendency to go broody. A silkie hen can produce 100 eggs in an ideal year. Their capacity for incubation, which has been selectively bred out of most fowl bred especially for egg production, is often exploited by poultry keepers by allowing Silkies to raise the offspring of other birds.
In cuisine
The black meat of a Silkie is generally considered an unusual attribute in European and American cuisines. In contrast, several Asian cuisines consider Silkie meat a gourmet food. Chinese cuisine especially values the breed, but it is also a common ingredient in some Japanese, Cambodian, Vietnamese and Korean dishes. Areas where Chinese cuisine has a strong influence, such as Malaysia, may also cook Silkie. As early as the 7th century, traditional Chinese medicine has held that chicken soup made with Silkie meat is a curative food. The usual methods of cooking include using Silkie to make broth, braising, and in curries. Traditional Chinese soup made with Silkie also uses ingredients such as wolfberries, Dioscorea polystachya (mountain yam), aged dried citrus peel, and fresh ginger. A few fusion restaurants in metropolitan areas of the West have also cooked it as a part of traditional American or French cuisine, such as in confit.
| Biology and health sciences | Chickens | Animals |
431003 | https://en.wikipedia.org/wiki/Zucchini | Zucchini | The zucchini (; : zucchini or zucchinis), courgette () or baby marrow (Cucurbita pepo) is a summer squash, a vining herbaceous plant whose fruit are harvested when their immature seeds and epicarp (rind) are still soft and edible. It is closely related, but not identical, to the marrow; its fruit may be called marrow when mature.
Ordinary zucchini fruit are any shade of green, though the golden zucchini is a deep yellow or orange. At maturity, they can grow to nearly in length, but they are normally harvested at about .
In botany, the zucchini's fruit is a pepo, a berry (the swollen ovary of the zucchini flower) with a hardened epicarp. In cookery, it is treated as a vegetable, usually cooked and eaten as an accompaniment or savory dish, though occasionally used in sweeter cooking.
Zucchini occasionally contain toxic cucurbitacins, making them extremely bitter, and causing severe gastero-enteric upsets. Causes include stressed growing conditions, and cross pollination with ornamental squashes.
Zucchini descends from squashes first domesticated in Mesoamerica over 7,000 years ago, but the zucchini itself was bred in Milan in the late 19th century.
Naming and etymology
The plant has three names in English, all of them meaning 'small marrow': zucchini (an Italian loanword), usually used in the plural form even when only one zucchino is meant, courgette (a French loanword), and baby marrow (South African English). Zucchini and courgette are doublets, both descending from the Latin .
Zucchini
The name zucchini is used in American, Australian, Canadian and New Zealand English. It is loaned from Italian, where is the plural masculine diminutive of ().
In Italian, the masculine (plural: ) is attested earlier and hence preferred by the Accademia della Crusca, the Italian language regulator. The feminine (plural: ) is also found, and preferred by the Italian-language encyclopedia Treccani, which considers to be a Tuscan dialect word.
Zucchini is also used in Canadian French, Danish, German, and Swedish.
Courgette
The name courgette is used in British, Hiberno-, Malaysian, New Zealand, and South African English. It is loaned from French, where () is a diminutive of .
Courgette is also used in Dutch.
Baby marrow
The name baby marrow is used interchangeably in South Africa with courgette.
Flower
The female flower is a golden blossom on the end of each emergent zucchini. The male flower grows directly on the stem of the zucchini plant in the leaf axils (where leaf petiole meets stem), on a long stalk, and is slightly smaller than the female. Both flowers are edible and are often used to dress a meal or to garnish the cooked fruit.
Firm and fresh blossoms that are only slightly open are cooked to be eaten, with pistils removed from female flowers, and stamens removed from male flowers. The stems on the flowers can be retained as a way of giving the cook something to hold onto during cooking, rather than injuring the delicate petals, or they can be removed prior to cooking, or prior to serving. There are a variety of recipes in which the flowers may be deep fried as fritters or tempura (after dipping in a light tempura batter), stuffed, sautéed, baked, or used in soups.
History
Zucchini, like all squash, has its ancestry in the Americas, specifically Mesoamerica. The varieties of green, cylindrical squash harvested immature and typically called "zucchini" were cultivated in northern Italy, as much as three centuries after the introduction of cucurbits from the Americas. It appears that this occurred in the second half of the 19th century, although the first description of the variety under the name zucchini occurs in a work published in Milan in 1901. Early varieties usually appended the names of nearby cities in their names.
The first records of zucchini in the United States date to the early 1920s. It was almost certainly taken to America by Italian immigrants and probably was first cultivated in the United States in California. A 1928 report on vegetables grown in New York State treats 'Zucchini' as one among 60 cultivated varieties of C. pepo.
Culinary uses
When used for food, zucchini are usually picked when under in length, when the seeds are still soft and immature. Mature zucchini can be long or more. These larger ones often have mature seeds and hard skins, requiring peeling and seeding. A zucchini with the flowers attached is a sign of a truly fresh and immature fruit, and it is especially sought after for its sweeter flavor.
Zucchini is usually served cooked. It can be prepared using a variety of cooking techniques, including steamed, boiled, grilled, stuffed and baked, barbecued, fried, or incorporated in other recipes such as soufflés. Raw grated zucchini can also be combined with flour and spices in a zucchini bread, similar to banana bread, or incorporated into a cake mix to make zucchini cake, similar to carrot cake. Its flowers can be eaten stuffed and are a delicacy when deep-fat-fried (e.g., tempura).
Zucchini has a delicate flavor and can be found simply cooked with butter or olive oil and herbs, or in more complex dishes. The skin is usually left in place. When frying zucchini, it is recommended to pat down cut sections to make them drier, similarly to what may be done with eggplant, in order to keep the slices' shape while cooking. Zucchini can also be eaten raw, sliced or shredded, in a cold salad, as well as lightly cooked in hot salads, as in Thai or Vietnamese recipes. Mature (larger-sized) zucchini are well-suited for cooking in breads.
Zucchinis can be cut with a spiralizer into noodle-like spirals and used as a low-carbohydrate substitute for pasta or noodles, often referred to as 'zoodles'.
In Australia, a popular dish is a frittata-like dish called zucchini slice.
In Bulgaria, zucchini may be fried and then served with a dip, made from yogurt, garlic, and dill. Another popular dish is oven-baked zucchini—sliced or grated—covered with a mixture of eggs, yogurt, flour, and dill.
In Egypt, zucchini may be cooked with tomato sauce, garlic, and onions.
In France, zucchini is a key ingredient in ratatouille, a stew of summer vegetable-fruits and vegetables prepared in olive oil and cooked for an extended time over low heat. The dish, originating near present-day Nice, is served as a side dish or on its own at lunch with bread. Zucchini may be stuffed with meat or with other fruits such as tomatoes or bell peppers in a dish called courgette farcie (stuffed zucchini).
In Greece, zucchini is usually fried, stewed or boiled with other fruits (often green chili peppers and eggplants). It is served as an hors d'œuvre or as a main dish, especially during fasting seasons. Zucchini is also stuffed with minced meat, rice, and herbs and served with avgolemono sauce. In several parts of Greece, the flowers of the plant are stuffed with white cheese, usually feta or mizithra, or with a mixture of rice, herbs, and occasionally minced meat. They are then deep-fried or baked in the oven with tomato sauce.
In Italy, zucchini is served in a variety of ways: fried, baked, boiled, or deep fried, alone or in combination with other ingredients. At home and in some restaurants, it is possible to eat the flowers, as well, deep-fried, known as fiori di zucca (cf. pumpkin flower fritter).
In the cuisines of the former Ottoman Empire, zucchini is often stuffed and called dolma. It is also used in various stews, both with and without meat, including ladera.
In Sephardic Jewish cuisine, medias (from Judeo-Spanish, meaning "halves") is a dish of halved zucchinis stuffed with meat and a mixture of ingredients, and cooked in a sour lemon sauce.
In Mexico, the flower (known as flor de calabaza) is often cooked in soups or used as a filling for quesadillas. The fruit is used in stews, soups (i.e. caldo de res, de pollo, or de pescado, mole de olla, etc.) and other preparations. The flower, as well as the fruit, is eaten often throughout Latin America.
In Russia, Ukraine and other CIS countries, zucchini usually is coated in flour or semolina and then fried or baked in vegetable oil, served with sour cream. Another popular recipe is "zucchini caviar", a squash spread made from thermically processed zucchini, carrots, onions and tomato paste, produced either at home or industrially as a vegetable preserve.
In Turkey, zucchini is the main ingredient in the popular dish mücver, or "zucchini pancakes", made from shredded zucchini, flour, and eggs, lightly fried in olive oil and eaten with yogurt. They are also often used in kebabs along with various meats. The flowers are also used in a cold dish, where they are stuffed with a rice mix with various spices and nuts and stewed.
In the United States, fried zucchini was invented in Pittsburgh.
In 2005, a poll of 2,000 people revealed it to be Britain's 10th favorite culinary vegetable.
Stuffed zucchini is found in many cuisines. As an example, in Lebanon, zucchini can be used to create Kousa Mahshi, which translates to "stuffed zucchini" in Arabic. The dish is made by coring the squash and then stuffing it with rice and spiced ground beef. Vegetables and other protein substitutes such as lamb may also be used. The contents of the zucchini are cooked by first boiling it and then reducing the heat of the zucchini's pot or container before letting it simmer for an hour. There's also Lebanese Zucchini Stew, or Mnazelah, a stew consisting of zucchini, potatoes, tomatoes, meat, and varied spices. Typical stuffings in the Middle Eastern family of dolma include rice, onions, tomato, and sometimes meat.
Nutrition
Zucchini are low in food energy (approximately per fresh zucchini) and contain good amounts of folate (24 μg/100 g), potassium (261 mg/100 g), provitamin A (200 IU [10 RAE]/100 g) and vitamin C (12.9 mg/100 g) .
Toxicology
Members of the plant family Cucurbitaceae, which includes zucchini / marrows, pumpkins and cucumbers, can contain toxins called cucurbitacins. These are steroids which defend the plants from predators, and have a bitter taste to humans. Cultivated cucurbitaceae are bred for low levels of the toxin and are safe to eat. However, ornamental pumpkins can have high levels of cucurbitacins, and such ornamental plants can cross-fertilize edible cucurbitaceae—any such cross-fertilized seeds used by the gardener for growing food in the following season can therefore potentially produce bitter and toxic fruit. Dry weather or irregular watering can also favor the production of the toxin, which is not destroyed by cooking. Humans with an impaired sense of taste (particularly the elderly) should therefore ask another person to taste the zucchini for them. This toxin has caused at least one death of an elderly person, in 2015. Investigators warned that gardeners should not save their own seeds, as reversion to forms containing more poisonous cucurbitacin might occur.
Zucchini can also be responsible for allergy caused by the presence of a protein: profilin. The sap released when peeling young zucchini also contains a viscous substance which when drying on the hands gives the impression of super-glue and dry hands.
Cultivation
Zucchini is very easy to cultivate in temperate climates. As such, it has a reputation among home gardeners for overwhelming production. The part harvested as "zucchini" is the immature fruit, although the flowers, mature fruit, and leaves are eaten, as well. One good way to control overabundance is to harvest the flowers, which are an expensive delicacy in markets because of the difficulty in storing and transporting them. The male flower is borne on the end of a stalk and is longer-lived.
While easy to grow, zucchini, like all squash, requires plentiful bees for pollination. In areas of pollinator decline or high pesticide use, such as mosquito-spray districts, gardeners often experience fruit abortion, where the fruit begins to grow, then dries or rots. This is due to an insufficient number of pollen grains delivered to the female flower. It can be corrected by hand pollination or by increasing the bee population.
Closely related to zucchini are Lebanese summer squash or kusa (not to be confused with cushaw), but they often are lighter green or even white. Some seed catalogs do not distinguish them. Various varieties of round zucchinis are grown in different countries under different names, such as "Tondo di Piacenza" in Italy, "Qarabaghli" in Malta and "Ronde de Nice" in France. In the late 1990s, American producers in California cultivated and began marketing round yellow and green zucchini known as "8-ball" squash (the yellow ones are sometimes known as "1-ball" or "gold ball"). White zucchini (summer squash) is sometimes seen as a mutation and can appear on the same plant as its green counterpart.
Cultivars
Bianco di Trieste
Black Beauty, very dark green
Cocozelle, dark green with white stripes, heirloom
| Biology and health sciences | Cucurbitales | null |
431109 | https://en.wikipedia.org/wiki/Petrified%20wood | Petrified wood | Petrified wood (from Ancient Greek meaning 'rock' or 'stone'; literally 'wood turned into stone'), is the name given to a special type of fossilized wood, the fossilized remains of terrestrial vegetation. Petrifaction is the result of a tree or tree-like plants having been replaced by stone via a mineralization process that often includes permineralization and replacement. The organic materials making up cell walls have been replicated with minerals (mostly silica in the form of opal, chalcedony, or quartz). In some instances, the original structure of the stem tissue may be partially retained. Unlike other plant fossils, which are typically impressions or compressions, petrified wood is a three-dimensional representation of the original organic material.
The petrifaction process occurs underground, when wood becomes buried in water or volcanic ash. The presence of water reduces the availability of oxygen which inhibits aerobic decomposition by bacteria and fungi. Mineral-laden water flowing through the sediments may lead to permineralization, which occurs when minerals precipitate out of solution filling the interiors of cells and other empty spaces. During replacement, the plant's cell walls act as a template for mineralization. There needs to be a balance between the decay of cellulose and lignin and mineral templating for cellular detail to be preserved with fidelity. Most of the organic matter often decomposes, however some of the lignin may remain. Silica in the form of opal-A, can encrust and permeate wood relatively quickly in hot spring environments. However, petrified wood is most commonly associated with trees that were buried in fine grained sediments of deltas and floodplains or volcanic lahars and ash beds. A forest where such material has petrified becomes known as a petrified forest.
Formation
Petrified wood forms when woody stems of plants are buried in wet sediments saturated with dissolved minerals. The lack of oxygen slows decay of the wood, allowing minerals to replace cell walls and to fill void spaces in the wood.
Wood is composed mostly of holocellulose (cellulose and hemicellulose) and lignin. Together, these substances make up 95% of the dry composition of wood. Almost half of this is cellulose, which gives wood much of its strength. Cellulose is composed of long chains of polymerized glucose arranged into microfibrils that reinforce the cell walls in the wood. Hemicellulose, a branched polymer of various simple sugars, makes up the majority of the remaining composition of hardwood while lignin, which is a polymer of phenylpropanes, is more abundant in softwood. The hemicellulose and lignin encrust and reinforce the cellulose microfibrils.
Dead wood is normally rapidly decomposed by microorganisms, beginning with the holocellulose. The lignin is hydrophobic (water-repelling) and much slower to decay. The rate of decay is affected by temperature and moisture content, but exclusion of oxygen is the most important factor preserving wood tissue: Organisms that decompose lignin must have oxygen for their life processes. As a result, fossil wood older than Eocene (about 56 million years old or older) has lost almost all its holocellulose, and only lignin remains. In addition to microbial decomposition, wood buried in an alkaline environment is rapidly broken down by inorganic reactions with the alkali.
Wood is preserved from decomposition by rapid entombment in mud, particularly mud formed from volcanic ash. The wood is then mineralized to transform it to stone. Non-mineralized wood has been recovered from Paleozoic formations, particularly Callixylon from Berea Sandstone, but this is very unusual. The petrified wood is later exposed by erosion of surrounding sediments. Non-mineralized fossil wood is rapidly destroyed when exposed by erosion, but petrified wood is quite durable.
Some 40 minerals have been identified in petrified wood, but silica minerals are by far the most important. Calcite and pyrite are much less common, and others are quite rare. Silica binds to the cellulose in cell walls via hydrogen bonding and forms a kind of template. Additional silica then replaces the cellulose as it decomposes, so that cell walls are often preserved in great detail. Thus silicification begins within the cell walls, and the spaces within and between cells are filled with silica more gradually. Over time, almost all the original organic material is lost; only around 10% remains in the petrified wood. The remaining material is nearly pure silica, with only iron, aluminum, and alkali and alkaline earth elements present in more than trace amounts. Iron, calcium, aluminum are the most common, and one or more of these elements may make up more than 1% of the composition.
Just what form the silica initially takes is still a topic of research. There is evidence of initial deposition as opal, which then crystallizes to quartz over long time periods. On the other hand, there is some evidence that silica is deposited directly as quartz.
Wood can become silicified very rapidly in silica-rich hot springs. While wood petrified in this setting is only a minor part of the geologic record, hot spring deposits are important to paleontologists because such deposits sometimes preserve more delicate plant parts in exquisite detail. These Lagerstätte deposits include the Paleozoic Rhynie Chert and East Kirkton Limestone beds, which record early stages in the evolution of land plants.
Most of the color in petrified wood comes from trace metals. Of these, iron is the most important, and it can produce a range of hues depending on its oxidation state. Chromium produces bright green petrified wood. Variations in color likely reflect different episodes of mineralization. In some cases, variations may come from chromatographic separation of trace metals.
Wood can also be petrified by calcite, as occurs in concretions in coal beds. Wood petrified by calcite tends to retain more of its original organic material. Petrification begins with deposition of goethite in the cell walls, followed by deposition of calcite in the void spaces. Carbonized wood is resistant to silicification and is usually petrified by other minerals. Wood petrified by minerals other than silica minerals tends to accumulate heavy metals, such as uranium, selenium, and germanium, with uranium most common in wood high in lignin and germanium most common in wood preserved in coal beds. Boron, zinc, and phosphorus are anomalously low in fossil wood, suggesting they are leached away or scavenged by microorganisms.
Less commonly, the replacement minerals in petrified wood are chalcocite or other sulfide minerals. These have been mined as copper ore at locations such as the Nacimiento Mine near Cuba, New Mexico.
Simulated petrified wood
Scientists have attempted to duplicate the process of petrification of wood, both to better understand the natural petrification process and for its possible use as a ceramic material. Early attempts used sodium metasilicate as a source of silica, but tetraethyl orthosilicate has proven more promising.
Uses
Petrified wood has limited use in jewelry, but is mostly used for decorative pieces such as book ends, table tops, clock faces, or other ornamental objects. A number of Ancestral Puebloan structures near Petrified Forest National Park were constructed of petrified wood, including the Agate House Pueblo. Petrified wood is also used in New Age healing.
Occurrences
Petrified wood is found worldwide in sedimentary beds ranging in age from the Devonian (about 390 million years ago), when woody plants first appeared on dry land, to nearly the present. Petrified "forests" tend to be either entire ecosystems buried by volcanic eruptions, in which trunks often remain in their growth positions, or accumulations of drift wood in fluvial environments. Amethyst Ridge at Yellowstone National Park shows 27 successive forest ecosystems buried by eruptions, while Petrified Forest National Park is a particularly fine example of fluvial accumulations of driftwood.
Volcanic ash is particularly suitable for preservation of wood, because large quantities of silica are released as the ash weathers. The presence of petrified wood in a sedimentary bed is often an indication of the presence of weathered volcanic ash. Petrified wood can also form in arkosic sediments, rich in feldspar and other minerals that release silica as they break down. The warm super monsoon climates of the Carboniferous through Permian periods seem to have favored this process. Preservation of petrified forests in volcanic ash beds is less affected by climate and preserves a greater diversity of species.
Areas with a large number of petrified trees include:
Africa
Egypt – petrified forest in Cairo-Suez road, declared a national protectorate by the ministry of environment, also in the area of New Cairo at the Extension of Nasr City, El Qattamiyya, near El Maadi district, and Al Farafra oasis.
Libya – Great Sand Sea – Hundreds of square miles of petrified trunks, branches and other debris mixed with Stone Age artifacts
Madagascar – Northwest Coast
Namibia – petrified forest of Damaraland
Sudan – petrified forest north of El-Kurru
Asia
China – in the Junggar Basin of Xinjiang, northwest China, the government has issued a crackdown on collecting of this material.
India – protected geological sites known for petrified wood are the National Fossil Wood Park, Tiruvakkarai (20-million-year-old fossils), and the Akal Wood Fossil Park (180-million-year-old fossils). Petrified wood has also been discovered in Dholavira in Kutch, Gujarat, dating back to 187–176 million years.
Japan – there is a fossilized forest preserved at Sendai City Tomizawa Site Museum
Indonesia – petrified wood covers several areas in Banten and also in some part of Mount Halimun Salak National Park.
Israel – several examples of petrified wood occur in the HaMakhtesh HaGadol in the Negev desert.
Pakistan – Sindh – Dadu – Petrified Forest at Khirthar National Park
Saudi Arabia – petrified forest north of Riyadh
Thailand – Bantak Petrified Forest Park in Ban Tak District has the longest petrified log in the world, officially measuring 69.7 metres.
Oceania
Australia – has deposits of petrified and opalized wood. Chinchilla, Queensland is famous for its 'Chinchilla Red'.
New Zealand:
Curio Bay on The Catlins coast contains many petrified wood examples.
Fossil Forest, Takapuna, Auckland, New Zealand
Titahi Bay Beach, Porirua, New Zealand
Europe
Belgium – Geosite Goudberg near Hoegaarden.
Czech Republic, Nová Paka – The most famous locality on Permian–Carboniferous rocks in the Czech Republic.
France – petrified forest in the village of Champclauson
Georgia – Goderdzi Petrified Forest Natural Monument.
Germany – the museum of natural history in Chemnitz has a collection of petrified trees, from the in situ Chemnitz petrified forest, found in the town in 1737.
Greece – Petrified forest of Lesvos, at the western tip of the island of Lesbos, is possibly the largest of the petrified forests, covering an area of over and declared a National Monument in 1985. Large, upright trunks complete with root systems can be found, as well as trunks up to 22 m in length.
Italy:
, petrified forest near Avigliano Umbro, Umbria (Central Italy), age Piacenzian.
, petrified forest near Soddì (Province of Oristano, Sardinia), age Chattian–Aquitanian.
Norway – Fossilized tropical forest in Svalbard
Ukraine – petrified araucaria trunks near Druzhkivka
United Kingdom
Fossil Grove, Glasgow, Scotland
Fossil Forest, Dorset, England
North America
Canada – in the badlands of southern Alberta; petrified wood is the provincial stone of Alberta. Axel Heiberg Island in Nunavut has a large petrified forest. In and around the North Saskatchewan river, around the Edmonton area. Blanche Brook, in Stephenville, Newfoundland, has 305-million-year-old examples.
United States
Petrified Wood Park in Lemmon, South Dakota
Ginkgo/Wanapum State Park in Washington state
Petrified Forest National Park in Arizona
Petrified Forest in California
Mississippi Petrified Forest in Flora, Mississippi
Cherokee Ranch petrified forest
Florissant Fossil Beds National Monument near Florissant, Colorado
Yellowstone Petrified Forest and Gallatin Petrified Forest, Yellowstone National Park, Wyoming
The south unit of Theodore Roosevelt National Park outside Medora, North Dakota
Gilboa Fossil Forest in New York
Black Hills Petrified Forest in South Dakota
Escalante Petrified Forest State Park in Utah
Agate Desert in the Upper Rogue River Valley near Medford, Oregon
Fossil Forest in the Catskill region near Cairo, New York
Valley of Fire State Park in Nevada
Bisti Badlands in New Mexico
Prince William Forest Park in Virginia
South America
Argentina – the Sarmiento Petrified Forest and Jaramillo Petrified Forest in Santa Cruz Province in the Argentine Patagonia have many trees that measure more than in diameter and long.
Brazil:
in the geopark of Paleorrota, there is a vast area with petrified trees.
In the Heritage forest
('Fossil Trees Natural Monument') in Tocantins: petrified forests of dicksoniaceae (specifically Psaronius and Tietea) and arthropitys
Petrified forests of dicksoniaceae (specifically Psaronius and Tietea singularis) and arthropitys can also be found in the state of São Paulo
near Rio Poti, Piauí, Permian (around 280–270 million years ago).
Ecuador – Puyango Petrified Forest – One of the largest collections of petrified wood in the world.
| Biology and health sciences | Paleontology | Biology |
431226 | https://en.wikipedia.org/wiki/Sawfly | Sawfly | Sawflies are wasp-like insects that are in the suborder Symphyta within the order Hymenoptera, alongside ants, bees, and wasps. The common name comes from the saw-like appearance of the ovipositor, which the females use to cut into the plants where they lay their eggs. The name is associated especially with the Tenthredinoidea, by far the largest superfamily in the suborder, with about 7,000 known species; in the entire suborder, there are 8,000 described species in more than 800 genera. Symphyta is paraphyletic, consisting of several basal groups within the order Hymenoptera, each one rooted inside the previous group, ending with the Apocrita which are not sawflies.
The primary distinction between sawflies and the Apocrita – the ants, bees, and wasps – is that the adults lack a "wasp waist", and instead have a broad connection between the abdomen and the thorax. Some sawflies are Batesian mimics of wasps and bees, and the ovipositor can be mistaken for a stinger. Sawflies vary in length, most measuring ; the largest known sawfly measured . The larvae are caterpillar-like, but can be distinguished by the number of prolegs and the absence of crochets in sawfly larvae. The great majority of sawflies are plant-eating, though the members of the superfamily Orussoidea are parasitic.
Predators include birds, insects and small animals. The larvae of some species have anti-predator adaptations such as regurgitating irritating liquid and clustering together for safety in numbers. Sawflies are hosts to many parasitoids, most of which are Hymenoptera, the rest being Diptera.
Adult sawflies are short-lived, with a life expectancy of 7–9 days, though the larval stage can last from months to years, depending on the species. Parthenogenetic females, which do not need to mate to produce fertilised eggs, are common in the suborder, though many species have males. The adults feed on pollen, nectar, honeydew, sap, other insects, including hemolymph of the larvae hosts; they have mouth pieces adapted to these types of feeding.
Sawflies go through a complete metamorphosis with four distinct life stages – egg, larva, pupa and adult. The female uses her ovipositor to drill into plant material (or, in the case of Orussoidea, other insects) and then lays eggs in groups called rafts or pods. After hatching, larvae feed on plants, often in groups. As they approach adulthood, the larvae seek a protected spot to pupate, typically in bark or the soil. Large populations of species such as the pine sawfly can cause substantial damage to economic forestry, while others such as the iris sawfly are major pests in horticulture. Outbreaks of sawfly larvae can defoliate trees and may cause dieback, stunting or death. Sawflies can be controlled through the use of insecticides, natural predators and parasitoids, or mechanical methods.
Sawflies first appeared 250 million years ago in the Triassic. The oldest superfamily, the Xyeloidea, has existed into the presents. Over 200million years ago, a lineage of sawflies evolved a parasitoid lifestyle, with carnivorous larvae that ate the eggs or larvae of other insects. Sawflies are distributed globally, though they are more diverse in the northernmost hemispheres.
Etymology
The suborder name "Symphyta" derives from the Greek word , meaning 'grown together', referring to the group's distinctive lack of a wasp waist between prostomium and peristomium. Its common name, "sawfly", derives from the saw-like ovipositor that is used for egg-laying, in which a female makes a slit in either a stem or plant leaf to deposit the eggs. The first known use of this name was in 1773. Sawflies are also known as "wood-wasps".
Phylogeny
In his original description of Hymenoptera in 1863, German zoologist Carl Gerstaecker divided them into three groups, Hymenoptera aculeata, Hymenoptera apocrita and Hymenoptera phytophaga. However, four years later in 1867, he described just two groups, H. apocrita syn. genuina and H. symphyta syn. phytophaga. Consequently, the name Symphyta is given to Gerstaecker as the zoological authority. In his description, Gerstaecker distinguished the two groups by the transfer of the first abdominal segment to the thorax in the Apocrita, compared to the Symphyta. Consequently, there are only eight dorsal half segments in the Apocrita, against nine in the Symphyta. The larvae are distinguished in a similar way.
The Symphyta have therefore traditionally been considered, alongside the Apocrita, to form one of two suborders of Hymenoptera. Symphyta are the more primitive group, with comparatively complete venation, larvae that are largely phytophagous, and without a "wasp-waist", a symplesiomorphic feature. Together, the Symphyta make up less than 10% of hymenopteran species. While the terms sawfly and Symphyta have been used synonymously, the Symphyta have also been divided into three groups, true sawflies (phyllophaga), woodwasps or xylophaga (Siricidae), and Orussidae. The three groupings have been distinguished by the true sawflies' ventral serrated or saw-like ovipositor for sawing holes in vegetation to deposit eggs, while the woodwasp ovipositor penetrates wood and the Orussidae behave as external parasitoids of wood-boring beetles. The woodwasps themselves are a paraphyletic ancestral grade. Despite these limitations, the terms have utility and are common in the literature.
While most hymenopteran superfamilies are monophyletic, as is Hymenoptera, the Symphyta has long been seen to be paraphyletic. Cladistic methods and molecular phylogenetics are improving the understanding of relationships between the superfamilies, resulting in revisions at the level of superfamily and family. The Symphyta are the most primitive (basal) taxa within the Hymenoptera (some going back 250 million years), and one of the taxa within the Symphyta gave rise to the monophyletic suborder Apocrita (wasps, bees, and ants). In cladistic analyses the Orussoidea are consistently the sister group to the Apocrita.
The oldest unambiguous sawfly fossils date back to the Middle or Late Triassic. These fossils, from the family Xyelidae, are the oldest of all Hymenoptera. One fossil, Archexyela ipswichensis from Queensland is between 205.6 and 221.5million years of age, making it among the oldest of all sawfly fossils. More Xyelid fossils have been discovered from the Middle Jurassic and the Cretaceous, but the family was less diverse then than during the Mesozoic and Tertiary. The subfamily Xyelinae were plentiful during these time periods, in which Tertiary faunas were dominated by the tribe Xyelini; these are indicative of a humid and warm climate.
The cladogram is based on Schulmeister 2003.
Taxonomy
There are approximately 8,000 species of sawfly in more than 800 genera, although new species continue to be discovered. However, earlier studies indicated that 10,000 species grouped into about 1,000 genera were known. Early phylogenies such as that of Alexandr Rasnitsyn, based on morphology and behaviour, identified nine clades which did not reflect the historical superfamilies. Such classifications were replaced by those using molecular methods, starting with Dowton and Austin (1994). As of 2013, the Symphyta are treated as nine superfamilies (one extinct) and 25 families. Most sawflies belong to the Tenthredinoidea superfamily, with about 7,000 species worldwide. Tenthredinoidea has six families, of which Tenthredinidae is by far the largest with some 5,500 species.
Extinct taxa are indicated by a dagger ().
Superfamilies and families
Superfamily Anaxyeloidea
Family Anaxyelidae (1 species) and 12 genera
Superfamily Cephoidea (1 and 1 family)
Family Cephidae (21 genera, 160 spp. and 3 genera
Superfamily Karatavitoidea (1 family)
Family Karatavitidae (7 genera)
Superfamily Orussoidea (1 and 1 family)
Family Orussidae (16 genera, 82 spp.) and 3 genera
Superfamily Pamphilioidea (2 and 1 families) (syn. Megalodontoidea)
Family Megalodontesidae (1 genera, 42 spp.) and 1 genus
Family Pamphiliidae (10 genera, 291 spp.) and 3 genera
Superfamily Siricoidea (2 and 5 families)
Family Siricidae (11 genera, 111 spp.) and 9 genera
Superfamily Tenthredinoidea (6 and 2 families)
Family Argidae (58 genera, 897 spp.) and 1 genus
Family Blasticotomidae (2 genera, 12 spp.) and 1 genus
Family Cimbicidae (16 genera, 182 spp.) and 6 genera
Family Diprionidae (11 genera, 136 spp.) and 2 genera
Family Pergidae (60 genera, 442 spp.)
Family Tenthredinidae (400 genera, 5,500 spp.) and 14 genera
Superfamily Xiphydrioidea
Family Xiphydriidae (28 genera, 146 spp.)
Superfamily Xyeloidea
Family Xyelidae (5 genera, 63 spp.) and 47 genera
Description
Many species of sawfly have retained their ancestral attributes throughout time, specifically their plant-eating habits, wing veins and the unmodified abdomen, where the first two segments appear like the succeeding segments. The absence of the narrow wasp waist distinguishes sawflies from other members of hymenoptera, although some are Batesian mimics with similar to wasps and bees, and the ovipositor can be mistaken for a stinger. Most sawflies are stubby and soft-bodied, and fly weakly. Sawflies vary in length: Urocerus gigas, which can be mistaken as a wasp due to its black-and-yellow striped body, can grow up to in length, but among the largest sawflies ever discovered was Hoplitolyda duolunica from the Mesozoic, with a body length of and a wingspan of . The smaller species only reach lengths of .
Heads of sawflies vary in size, shape and sturdiness, as well as the positions of the eyes and antennae. They are characterised in four head types: open head, maxapontal head, closed head and genapontal head. The open head is simplistic, whereas all the other heads are derived. The head is also hypognathous, meaning that the lower mouthparts are directed downwards. When in use, the mouthparts may be directed forwards, but this is only caused when the sawfly swings its entire head forward in a pendulum motion. Unlike most primitive insects, the sutures (rigid joints between two or more hard elements on an organism) and sclerites (hardened body parts) are obsolescent or absent. The clypeus (a sclerite that makes up an insects "face") is not divided into a pre- and postclypeus, but rather separated from the front. The antennal sclerites are fused with the surrounding head capsule, but these are sometimes separated by a suture. The number of segments in the antennae vary from six in the Accorduleceridae to 30 or more in the Pamphiliidae. The compound eyes are large with a number of facets, and there are three ocelli between the dorsal portions of the compound eyes. The tentorium comprises the whole inner skeleton of the head.
Three segments make up the thorax: the mesothorax, metathorax and prothorax, as well as the exoskeletal plates that connect with these segments. The legs have spurs on their fourth segments, the tibiae. Sawflies have two pairs of translucent wings. The fore and hind wings are locked together with hooks. Parallel development in sawfly wings is most frequent in the anal veins. In all sawflies, 2A and 3A tend to fuse with the first anal vein. This occurs in several families including Argidae, Diprionidae and Cimbicidae.
The larvae of sawflies are easily mistaken for lepidopteran larvae (caterpillars). However, several morphological differences can distinguish the two: while both larvae share three pairs of thoracic legs and an apical pair of abdominal prolegs, lepidopteran caterpillars have four pairs of prolegs on abdominal segments 3–6 while sawfly larvae have five pairs of prolegs located on abdominal segments 2–6; crochets are present on lepidopteran larvae, whereas on sawfly larvae they are not; the prolegs of both larvae gradually disappear by the time they burrow into the ground, therefore making it difficult to distinguish the two; and sawfly larvae only have a single pair of minute eyes, whereas lepidopteran larvae have four to six eyes on each side of the head. Sawfly larvae behave like lepidopteran larvae, walking about and eating foliage. Some groups have larvae that are eyeless and almost legless; these larvae make tunnels in plant tissues including wood. Many species of sawfly larvae are strikingly coloured, exhibiting colour combinations such as black and white while others are black and yellow. This is a warning colouration because some larvae can secrete irritating fluids from glands located on their undersides.
Distribution
Sawflies are widely distributed throughout the world. The largest family, the Tenthredinidae, with some 5,000 species, are found on all continents except Antarctica, though they are most abundant and diverse in the temperate regions of the northern hemisphere; they are absent from New Zealand and there are few of them in Australia. The next largest family, the Argidae, with some 800 species, is also worldwide, but is most common in the tropics, especially in Africa, where they feed on woody and herbaceous angiosperms. Of the other families, the Blasticotomidae and Megalodontidae are Palearctic; the Xyelidae, Pamphilidae, Diprionidae, Cimbicidae, and Cephidae are Holarctic, while the Siricidae are mainly Holarctic with some tropical species. The parasitic Orussidae are found worldwide, mostly in tropical and subtropical regions. The wood-boring Xiphydriidae are worldwide, but most species live in the subtropical parts of Asia.
Behaviour and ecology
Sawflies are mostly herbivores, feeding on plants that have a high concentration of chemical defences. These insects are either resistant to the chemical substances, or they avoid areas of the plant that have high concentrations of chemicals. The larvae primarily feed in groups; they are folivores, eating plants and fruits on native trees and shrubs, though some are parasitic. However, this is not always the case; Monterey pine sawfly (Itycorsia) larvae are solitary web-spinners that feed on Monterey pine trees inside a silken web. The adults feed on pollen and nectar.
Sawflies are eaten by a wide variety of predators. While many birds find the larvae distasteful, some such as the currawong (Strepera) and stonechats (Saxicola) eat both adults and larvae. The larvae are an important food source for the chicks of several birds, including partridges. Sawfly and moth larvae form one third of the diet of nestling corn buntings (Emberiza calandra), with sawfly larvae being eaten more frequently on cool days. Black grouse (Tetrao tetrix) chicks show a strong preference for sawfly larvae. Sawfly larvae formed 43% of the diet of chestnut-backed chickadees (Poecile rufescens). Small carnivorous mammals such as the masked shrew (Sorex cinereus), the northern short-tailed shrew (Blarina brevicauda) and the deer mouse (Peromyscus maniculatus) predate heavily on sawfly cocoons. Insects such as ants and certain species of predatory wasps (Vespula vulgaris) eat adult sawflies and the larvae, as do lizards and frogs. Pardalotes, honeyeaters and fantails (Rhipidura) occasionally consume laid eggs, and several species of beetle larvae prey on the pupae.
The larvae have several anti-predator adaptations. While adults are unable to sting, the larvae of species such as the spitfire sawfly regurgitate a distasteful irritating liquid, which makes predators such as ants avoid the larvae. In some species, the larvae cluster together, reducing their chances of being killed, and in some cases form together with their heads pointing outwards or tap their abdomens up and down. Some adults bear black and yellow markings that mimic wasps.
Parasites
Sawflies are hosts to many parasitoids, most of which are parasitic Hymenoptera; more than 40 species are known to attack them. However, information regarding these species is minimal, and fewer than 10 of these species actually cause a significant impact on sawfly populations. Many of these species attack their hosts in the grass or in other parasitoids. Well known and important parasitoids include Braconidae, Eulophidae and Ichneumonidae. Braconid wasps attack sawflies in many regions throughout the world, in which they are ectoparasitoids, meaning that the larvae live and feed outside of the hosts body; braconids have more of an impact on sawfly populations in the New World than they do in the Old World, possibly because there are no ichneumonid parasitoids in North America. Some braconid wasps that attack sawflies include Bracon cephi, B. lisogaster, B. terabeila and Heteropilus cephi. Female braconids locate sawfly larvae through the vibrations they produce when feeding, followed by inserting the ovipostior and paralysing the larva before laying eggs inside the host. These eggs hatch inside the larva within a few days, where they feed on the host. The entire host's body may be consumed by the braconid larvae, except for the head capsule and epidermis. The larvae complete their development within two or three weeks.
Ten species of wasps in the family Ichneumonidae attack sawfly populations, although these species are usually rare. The most important parasitoids in this family are species in the genus Collyria. Unlike braconids, the larvae are endoparasitoids, meaning that the larvae live and feed inside the hosts body. One well known ichneumonid is Collyria coxator, which is a dominant parasitoid of C. pygmaeus. Recorded parasitism rates in Europe are between 20–76%, and as many as eight eggs can be found in a single larva, but only one Collyria individual will emerge from its host. The larva may remain inside of their host until spring, where it emerges and pupates.
Several species in the family Eulophidae attack sawflies, although their impact is low. Two species in the genus Pediobius have been studied; the two species are internal larval parasitoids and have only been found in the northern hemisphere. Parasitism of sawflies by eulophids in grass exceeds 50%, but only 5% in wheat. It is unknown as to why the attack rate in wheat is low. Furthermore, some fungal and bacterial diseases are known to infect eggs and pupa in warm wet weather.
Outbreaks of certain sawfly species, such as Diprion polytomum, have led scientists to investigate and possibly collect their natural enemies to control them. Parasites of D. polytomum have been extensively investigated, showing that 31 species of hymenopterous and dipterous parasites attack it. These parasites have been used in successful biological control against pest sawflies, including Cephus cinctus throughout the 1930s and 1950s and C. pygmaeus in the 1930s and 1940s.
Life cycle and reproduction
Like all other hymenopteran insects, sawflies go through a complete metamorphosis with four distinct life stages – egg, larva, pupa and adult. Many species are parthenogenetic, meaning that females do not need fertilization to create viable eggs. Unfertilized eggs develop as male, while fertilized eggs develop into females (arrhenotoky). The lifespan of an individual sawfly is two months to two years, though the adult life stage is often very short (approximately 7 – 9 days), only long enough for the females to lay their eggs. The female uses its ovipositor to drill into plant material to lay her eggs (though the family Orussoidea lay their eggs in other insects). Plant-eating sawflies most commonly are associated with leafy material but some specialize on wood, and the ovipositors of these species (such as the family Siricidae) are specially adapted for the task of drilling through bark. Once the incision has been made, the female will lay as many as 30 to 90 eggs. Females avoid the shade when laying their eggs because the larvae develop much slower and may not even survive, and they may not also survive if they are laid on immature and glaucous leaves. Hence, female sawflies search for young adult leaves to lay their eggs on.
These eggs hatch in two to eight weeks, but such duration varies by species and also by temperature. Until the eggs have hatched, some species such as the small brown sawfly will remain with them and protects the eggs by buzzing loudly and beating her wings to deter predators. There are six larval stages that sawflies go through, lasting 2 – 4 months, but this also depends on the species. When fully grown, the larvae emerge from the trees en masse and burrow themselves into the soil to pupate. During their time outside, the larvae may link up to form a large colony if many other individuals are present. They gather in large groups during the day which gives them protection from potential enemies, and during the night they disperse to feed. The emergence of adults takes awhile, with some emerging anywhere between a couple months to 2 years. Some will reach the ground to form pupal chambers, but others may spin a cocoon attached to a leaf. Larvae that feed on wood will pupate in the tunnels they have constructed. In one species, the jumping-disc sawfly (Phyllotoma aceris) forms a cocoon which can act like a parachute. The larvae live in sycamore trees and do not damage the upper or lower cuticles of leaves that they feed on. When fully developed, they cut small perforations in the upper cuticle to form a circle. After this, they weave a silk hammocks within the circle; this silk hammock never touches the lower cuticle. Once inside, the upper-cuticle's disc separates and descends towards the surface with the larvae attaching themselves to the hammock. Once they reach the round, the larvae work their way into a sheltered area by jerking their discs along.
The majority of sawfly species produce a single generation per year, but others may only have one generation every two years. Most sawflies are also female, making males rare.
Relationship with humans
Sawflies are major economic pests of forestry. Species in the Diprionidae, such as the pine sawflies, Diprion pini and Neodiprion sertifer, cause serious damage to pines in regions such as Scandinavia. D. pini larvae defoliated in the largest outbreak in Finland, between 1998 and 2001. Up to 75% of the trees may die after such outbreaks, as D. pini can remove all the leaves late in the growing season, leaving the trees too weak to survive the winter. Little damage to trees only occurs when the tree is large or when there is minimal presence of larvae. Eucalyptus trees can regenerate quickly from damage inflicted by the larvae; however, they can be substantially damaged from outbreaks, especially if they are young. The trees can be defoliated completely and may cause "dieback", stunting or even death.
Sawflies are serious pests in horticulture. Different species prefer different host plants, often being specific to a family or genus of hosts. For example, Iris sawfly larvae, emerging in summer, can quickly defoliate species of Iris including the yellow flag and other freshwater species. Similarly the rose sawflies, Arge pagana and A. ochropus, defoliate rose bushes.
The giant woodwasp or horntail, Urocerus gigas, has a long ovipositor, which with its black and yellow colouration make it a good mimic of a hornet. Despite the alarming appearance, the insect cannot sting. The eggs are laid in the wood of conifers such as Douglas fir, pine, spruce, and larch. The larvae eat tunnels in the wood, causing economic damage.
Alternative measures to control sawflies can be taken. Small-scale, mechanical methods include visually confirming larval presence on a plant and subsequently removing them, either by pruning damaged leaves or removing the larvae from the leaves they are on. Larvae typically try to remain hidden on the underside of foliage. Upon removing larvae and/or the affected leaves from plants, they may be dispatched by squishing, or, alternatively, the cut leaves with larvae still attached may be fed to birds; if larger animals do not prey upon them, other insects will. However, this is not practical or useful for some, thus the larvae can be quickly dispatched by simply dropping foliage into a vessel of plain or saltwater, diluted hydrogen peroxide or isopropyl alcohol, insecticidal soap, or other garden chemical. In large-scale, industrial settings, where beneficial insect predators can also be used to eliminate larvae, as well as parasites, which have both been previously used in control programs. Small trees can be sprayed with a number of chemicals, including maldison, dimethoate, carbaryl, imidacloprid, etc., if removing larvae from trees is not effective enough.
| Biology and health sciences | Hymenoptera | Animals |
431310 | https://en.wikipedia.org/wiki/Chemical%20engineer | Chemical engineer | A chemical engineer is a professional equipped with the knowledge of chemistry and other basic sciences who works principally in the chemical industry to convert basic raw materials into a variety of products and deals with the design and operation of plants and equipment.<ref>MobyDick Dictionary of Engineering", McGraw-Hill, 2nd Ed.</ref> This person applies the principles of chemical engineering in any of its various practical applications, such as
Design, manufacture, and operation of plants and machinery in industrial chemical and related processes ("chemical process engineers");
Development of new or adapted substances for products ranging from foods and beverages to cosmetics to cleaners to pharmaceutical ingredients, among many other products ("chemical product engineers");
Development of new technologies such as fuel cells, hydrogen power and nanotechnology, as well as working in fields wholly or partially derived from chemical engineering such as materials science, polymer engineering, and biomedical engineering. This can include working of geophysical projects such as rivers, stones, and signs.
History
The president of the Institution of Chemical Engineers said in his presidential address "I believe most of us would be willing to regard Edward Charles Howard (1774–1816) as the first chemical engineer of any eminence". Others have suggested Johann Rudolf Glauber (1604–1670) for his development of processes for the manufacture of the major industrial acids.
The term appeared in print in 1839, though from the context it suggests a person with mechanical engineering knowledge working in the chemical industry.
In 1880, George E. Davis wrote in a letter to Chemical News "A Chemical Engineer is a person who possesses chemical and mechanical knowledge, and who applies that knowledge to the utilisation, on a manufacturing scale, of chemical action." He proposed the name Society of Chemical Engineers, for what was in fact constituted as the Society of Chemical Industry. At the first General Meeting of the Society in 1882, some 15 of the 300 members described themselves as chemical engineers, but the Society's formation of a Chemical Engineering Group in 1918 attracted about 400 members.
In 1905 a publication called The Chemical Engineer was founded in the US, and in 1908 the American Institute of Chemical Engineers was established.
In 1924 the Institution of Chemical Engineers adopted the following definition: "A chemical engineer is a professional man experienced in the design, construction and operation of plant and works in which matter undergoes a change of state and composition."
As can be seen from the later definition, the occupation is not limited to the chemical industry, but more generally the process industries, or other situations in which complex physical and/or chemical processes are to be managed.
The UK journal The Chemical Engineer'' (began 1956) has a series of biographies available online entitled “Chemical Engineers who Changed the World”,
Overview
Historically, the chemical engineer has been primarily concerned with process engineering, which can generally be divided into two complementary areas: chemical reaction engineering and separation processes. The modern discipline of chemical engineering, however, encompasses much more than just process engineering. Chemical engineers are now engaged in the development and production of a diverse range of products, as well as in commodity and specialty chemicals. These products include high-performance materials needed for aerospace, automotive, biomedical, electronic, environmental and military applications. Examples include ultra-strong fibers, fabrics, adhesives and composites for vehicles, bio-compatible materials for implants and prosthetics, gels for medical applications, pharmaceuticals, and films with special dielectric, optical or spectroscopic properties for opto-electronic devices. Additionally, chemical engineering is often intertwined with biology and biomedical engineering. Many chemical engineers work on biological projects such as understanding biopolymers (proteins) and mapping the human genome.
Employments and salaries
According to a 2015 salary survey by the American Institute of Chemical Engineers, the median annual salary for a chemical engineer was approximately $127,000. The survey was repeated in 2017 and the median annual salary dropped slightly to $124,000. The decrease in median salary was unexpected. A factor contributing to the decline may be that 2017’s survey was conducted by a different research and analysis firm. Median salaries ranged from $70,450 for chemical engineers with fewer than three years of experience to $156,000 for those with more than 40 years in the workforce.
In the UK, the IChemE 2016 Salary Survey reported a median salary of approximately £57,000, with a starting salary for a graduate averaging £28,350. Chemical engineering in the USA is one of the engineering disciplines with the highest participation of women, with 35% of students compared with 20% in engineering. In the UK in 2014, students starting degrees were 25% female, compared with 15% in engineering. US graduates who responded to a 2015 salary survey were 18.8% female.
According to the latest 2023 figures, Bayes Business School graduates get an average of £51,921 within 5 years of graduation, which is the most among UK universities. This was followed by the University of Oxford at £49,086 and the University of Warwick at £47,446.
| Physical sciences | Basics: General | Chemistry |
431462 | https://en.wikipedia.org/wiki/Water%20cannon | Water cannon | A water cannon is a device that shoots a high-velocity stream of water. Typically, a water cannon can deliver a large volume of water, often over dozens of meters. They are used in firefighting, large vehicle washing, riot control, and mining. Most water cannons fall under the category of a fire monitor.
Firefighting
Water cannons were first devised for use on fireboats. Extinguishing fires on boats and buildings near the water was much more difficult and dangerous before fireboats were invented. The first fireboat deployed in Los Angeles was commissioned on 1 August 1919. The first fireboat in New York City was Marine 1, deployed 1 February 1891. There may have been other fireboats elsewhere even earlier.
Fire trucks deliver water with much the same force and volume as water cannons, and have even been used in riot control situations, but are rarely referred to as water cannons outside this context.
Riot control
The first truck-mounted water cannon was used for riot control in Germany in the beginning of the 1930s.
The most modern versions do not expose the operator to the riot, and are controlled remotely from within the vehicle by a joystick. The Austrian-built WaWe 10.000 by Rosenbauer used by German police can carry of water, which can deploy water in all directions via three cannons, all of which are remotely controlled from inside the vehicle by a joystick. The vehicle has two forward cannons with a delivery rate of , and one rear cannon with a delivery rate of
Water cannons designed for riot control are still made in the United States and the United Kingdom, but most products are exported, particularly to Africa and parts of Asia such as Indonesia.
Safety
Use of water cannon in riot control contexts can lead to injury or death, with fatalities recorded in Indonesia (in 1996, when the cannon's payload contained ammonia), Zimbabwe (in 2007, when the use of cannons on a peaceful crowd caused panic), Turkey (in 2013, when the payload was laced with "liquid tear gas"), Ukraine (in 2014, with the death of activist and businessman Bogdan Kalynyak, reportedly catching pneumonia after being sprayed by a water cannon in freezing temperatures) and South Korea (in 2016, when a 68-year-old farmer died after injuries sustained by a water cannon the previous year).
Water cannons in use during the 1960s, which were generally adapted fire trucks, would knock protesters down and on occasion, tear their clothes.
On 30 September 2010, during a protest demonstration against the Stuttgart 21 project in Germany, a demonstrator was hit in the face by a water cannon. Dietrich Wagner, a retired engineer, suffered damage to his eyelids and retinas, resulting in near-complete loss of his eyesight. Graphic imagery was recorded of the event, sparking a national debate about police brutality and proportionality in the use of state force.
According to a report issued in the United Kingdom, using plastic bullets instead of water cannons was justified because the latter "are inflexible and indiscriminate", although several people had previously been killed or seriously injured by plastic bullets.
Media effect
The presence of the media at riots has had a significant impact on water cannon use. There is much pressure on police departments to avoid bad publicity, and water cannons often play badly in the press. It is considered that this is a likely reason that they are not used more often in some countries.
Confrontations that took place in the era of the American Civil Rights Movement, where water cannons were used by authorities to disperse crowds of protesting African Americans, has led to the demise of water cannons in the United States.
Alternative payload
Dye
In 1997 pink dye was reportedly added to the water used by South Korean and Indonesian police to disperse a riot. The implication is that they might use this mark to make it easier to arrest rioters later. The United Kingdom, which had sold the water cannon to Indonesia, condemned this practice (although the Royal Ulster Constabulary had used a water cannon with purple dye during The Troubles in Northern Ireland) but later approved the sale of further water cannons to them. Most modern water cannons are also capable of adding tear gas to the stream.
Electrified water jet
In 2004 Jaycor Tactical Systems was experimenting with additives (salt and additives to reduce the breakup of the stream into droplets) that would allow electricity to be conducted through water. They have demonstrated delivery from a distance of up to , but have not yet tested the device on people.
Although referred to as an electrified water cannon, this experiment involved a water jet much less powerful than a water cannon.
Other types
Water cannon differ from other similar devices in the volume of water delivered in a given time, the nozzle speed, the pressure that it is delivered at, and to a lesser extent the total volume that can be delivered. They are also generally portable. The method of employment is also important in labeling a device a water cannon. Nevertheless, the distinction between a water cannon and other similar devices is fuzzy. For example:-
Pressure washers generally produce an extremely high pressure stream where the power of the stream drops off significantly over a very short distance.
Water pistols and other toys deliver much lower volumes of water at a much lower pressure.
Ultra high pressure water jet cutters are used to cut a wide variety of materials including granite, concrete (see hydrodemolition), ceramics, fabric and even Kevlar. One such cutter delivers through a nozzle in diameter at 1 kilometre per second, which can cut a person at a close range. There are reports of accidental deaths involving the industrial use of high-pressure water.
Usage
Water cannon are still in large scale use in Chile, Belgium, the Netherlands and other parts of the world.
Australia
The New South Wales Police Force purchased a water cannon in 2007 and had it deployed on standby during an APEC meeting in Sydney that year. It was the first purchase of a water cannon by a police service in Australia. However, it ended up not being used during the APEC meeting, and was never used during any instance of civil unrest. Eventually it was retired and converted to a water tanker for fire department use.
Germany
The annual riots on 1 May in Berlin, the Schanzenfest fair in Hamburg, which regularly ends in riots, or other demonstrations, are usually accompanied by water cannon, which support riot police. The most commonly used water cannon in Germany over years was the Wasserwerfer 9000. Since 2019, the only water cannon type used by riot police, which are around 50 units in total, is the Wasserwerfer 10000.
Hong Kong
Three truck-mounted water cannon, officially known as 'Specialised Crowd Management Vehicles', were purchased by Hong Kong Police from France in mid-2018. The truck chassis were provided by Mercedes-Benz and the water spray devices were also made by German firm Ziegler. The three water cannon cost HK$27 million to purchase, a sum that was criticised as overpriced. The vehicles were frequently used by police on participants and bystanders during the 2019–20 Hong Kong protests. Blue dye was often added to the water to allow police to identify protesters. Pepper spray solution was also often added as an irritant.
On 20 October 2019, police used a water cannon to target and shoot a small group of pedestrians standing outside Kowloon Mosque, in Tsim Sha Tsui, using blue-dyed water mixed with a pepper solution. A large number of Hong Kong residents spontaneously went to the scene to clean up, with the incident resulting in an increased sense of inclusiveness among the Hong Kong public toward the city's Muslim and other minorities.
Israel
Since the 1980s, Israel has been exporting water cannons to numerous countries around the world. Bet Alpha Technologies, a company owned by Kibbutz Bet Alpha, has sold water cannons to Russia, China, Turkey, the United States, Latvia, Zambia, Argentina and Swaziland amounting to millions of dollars in sales. The Israel Police have made extensive use of water cannons during demonstrations. Its water canons are capable of spraying jets of water, paint (used to mark protesters for later arrest), gas, and Skunk in long or short pulses in an effective range of 40 meters. They are controlled controlled by a joystick and set of cameras and is equipped with a mine plow allows the vehicle to break through and push through hard barriers like barricades placed on the road. During the 2023 Israeli judicial reform protests, the Israel Police allegedly violated its own procedures when on several occasions they fired water streams directly toward protesters' heads, causing damage to the vision of some of them.
Thailand
During the 2020 Thai protests, on 16 October 2020, the police used water cannon claimed to have water containing an irritant that made protesters' eyes sting to disperse a peaceful protest in Bangkok.
Turkey
The Turkish police water cannon TOMA has been used against protesters many times, including the 2013 protests in Turkey, and are often present at protests of all sizes.
United Kingdom
Only six water cannons are operational in the United Kingdom, all held by the Police Service of Northern Ireland (PSNI); these are Somati RCV9000 Vehicle Mounted Water Cannons built on GINAF chassis, which after extensive evaluation by a Defence Scientific Advisory Council sub-committee as a less-lethal replacement of baton rounds, began to enter service with the PSNI from 2004 onwards. Water cannon use outside Northern Ireland is not approved, and would require the statutory authorisation from the Home Secretary for use in England and Wales or the parliament of Scotland for use in Scotland.
In June 2014, London's Deputy Mayor for Policing and Crime Stephen Greenhalgh authorised the Metropolitan Police to buy three-second-hand Wasserwerfer 9000s from the German Federal Police. Mayor of London Boris Johnson said that the purchase had been authorised before Parliamentary approval, as the three cannons cost £218,000 to purchase and would require a further £125,000 of work before being deemed suitable for service, as opposed to £870,000 for a single new machine. But after a study of their safety and effectiveness, Home Secretary Theresa May said in Parliament in July 2015 that she had decided not to license them for use. They were sold in November 2018 with the intention of them being dismantled for spare parts. The resale resulted in a net loss of £300,000.
United States
Truck-based water cannon, and fire hoses used as improvised water cannons, were used widely in the United States during the 1960s for both riot control and suppressing peaceful civil rights marches, including the infamous use ordered by Eugene "Bull" Connor in Birmingham, Alabama in 1963. The newsreel footage of police turning water cannons and police dogs on civilians—both student protesters and bystanders alike, including children as young as six—widely viewed as shocking and inappropriate and helped turn public sympathies towards civil rights. Water cannons were used in November 2016 during the Dakota Access Pipeline protests. In August 2020, state senator Floyd Prozanski suggested water cannons be used by police against protesters in Portland, Oregon.
The New York City Police Department previously had a water cannon made from a 1982 Oshkosh P-4 as part of their Disorder Control Unit, which was in their fleet until at least the 2000s. There are no recorded instances of it ever being deployed.
Mining
Water cannons are used in hydraulic mining to dislodge rock material or move sediment. In the placer mining of gold or tin, the resulting water-sediment slurry is directed through sluice boxes to remove the gold. It is also used in mining kaolin and coal.
Gallery
Other meanings
The term "water cannon" could also refer to:-
Similar land vehicles used for firefighting
Numerous large stationary toys
Waterjet in hydraulic mining
A type of railway wagon used to remove fallen leaves off the track: e.g. seen at Alexandra Palace on 25 October 2003
Tool for powerwashing large construction equipment.
| Technology | Less-lethal weapons | null |
431483 | https://en.wikipedia.org/wiki/Magpie%20goose | Magpie goose | The magpie goose (Anseranas semipalmata) is the sole living representative species of the family Anseranatidae. This common waterbird is found in northern Australia and southern New Guinea. As the species is prone to wandering, especially when not breeding, it is sometimes recorded outside its core range. The species was once also widespread in southern Australia but disappeared from there largely due to the drainage of the wetlands where the birds once bred. Due to their importance to Aboriginal people as a seasonal food source, as subjects of recreational hunting, and as a tourist attraction, their expansive and stable presence in northern Australia has been "ensured [by] protective management".
Description
Magpie geese are unmistakable birds with their black and white plumage and yellowish legs. The feet are only partially webbed, and the magpie goose feeds on vegetable matter in the water, as well as on land. Males are larger than females. Unlike true geese, their molt is gradual, so no flightless periods result. Their voice is a loud honking.
Systematics and evolution
This species is placed in the order Anseriformes, having the characteristic bill structure, but is considered to be distinct from the other species in this taxon. The related and extant families, Anhimidae (screamers) and Anatidae (ducks, geese, and swans), contain all the other taxa. The magpie goose is contained in the genus Anseranas and family Anseranatidae, which are monotypic now.
A cladistic study of the morphology of waterfowl found that the magpie goose was an early and distinctive offshoot, diverging after screamers and before all other ducks, geese, and swans.
This family is quite old, a living fossil, having apparently diverged before the Cretaceous–Paleogene extinction event — the relative Vegavis iaai lived some 68-67 million years ago. The fossil record is limited, nonetheless. The enigmatic genus Anatalavis (Hornerstown Late Cretaceous or Early Paleocene of New Jersey, USA - London Clay Early Eocene of Walton-on-the-Naze, England) is sometimes considered to be the earliest known. Other Paleogene birds sometimes considered magpie-geese are the genera Geranopsis from the Hordwell Formation Late Eocene to the Early Oligocene of England and Anserpica from the Late Oligocene of Billy-Créchy (France).
The earliest known member of the group in Australia is Eoanseranas represented by fossils found in the late Oligocene Carl Creek Limestone of Queensland. Additional fossils from North America and Europe suggest that the family was spread across the globe during the late Paleogene period. The Australian distribution of the living species ties in well with the presumed Gondwanan origin of Anseriformes, but Northern Hemisphere fossils are puzzling. Perhaps the magpie geese were one of the dominant groups of Paleogene waterfowl, only to become largely extinct later.
Ecology and status
The magpie goose is found in a variety of open wetland areas such as floodplains and swamps, where they wade and swim. They eat mostly vegetation such as dry grass blades, grass seeds, spike rush bulbs and wild rice.
Magpie geese are fairly sedentary apart from some movement during the dry season. They are colonial breeders and are gregarious outside of the breeding season when they can form large and noisy flocks of up to a few thousand individuals. Magpie geese nest on the ground or in trees where they can be five meters or higher above the ground. Their typical clutch is between 5-14 eggs. Some males mate with two females, all of which raise the young, unlike some other polygamous birds. This may be beneficial when predation of young is high as chicks raised by trios are more likely to survive.
This species is plentiful across its range, although this is significantly reduced in comparison to the range at time of European settlement. The range once extended as far south as the Coorong and the wetlands of the southeast of South Australia and Western Victoria. For Australia as a whole, it is not threatened and has a controlled hunting season when numbers are large. However, most of the southern populations were extirpated in the mid-20th century by overhunting and habitat destruction. The species has been subject to reintroduction projects such as Bool Lagoon between Penola and Naracoorte. Populations in more northern areas have again reached a level where it can be regularly utilized by hunters, although not in the example provided. The magpie goose was listed as near threatened on the 2007 advisory list of threatened vertebrate fauna in Victoria. In the December 2007 Flora and Fauna Guarantee Act list of threatened fauna, it is also listed. As of early 2008, an Action Statement for the recovery and future management of this species had not been prepared.
With the advent of climate change, and more frequent seawater inundations of the current extensive freshwater floodplains, CSIRO scientists argue that magpie geese populations may be at risk.
In Aboriginal languages
The Kunwinjku of western Arnhem Land know this bird as manimunak. It became an important food item with the formation of wetlands about 1500 ya, and is depicted in rock art from this period. Mimi figures are often shown holding goose-feather fans. In Yolŋu Matha the bird is known as gurrumaṯtji, or around Ramingining as gumang.
In the Wadawurrung language, the magpie goose is known as Ngangok.
Gallery
| Biology and health sciences | Anseriformes | Animals |
431529 | https://en.wikipedia.org/wiki/Supersolid | Supersolid | In condensed matter physics, a supersolid is a spatially ordered (i.e. solid) material with superfluid properties. In the case of helium-4, it has been conjectured since the 1960s that it might be possible to create a supersolid. Starting from 2017, a definitive proof for the existence of this state was provided by several experiments using atomic Bose–Einstein condensates. The general conditions required for supersolidity to emerge in a certain substance are a topic of ongoing research.
Background
A supersolid is a special quantum state of matter where particles form a rigid, spatially ordered structure, but also flow with zero viscosity. This is in contradiction to the intuition that flow, and in particular superfluid flow with zero viscosity, is a property exclusive to the fluid state, e.g., superconducting electron and neutron fluids, gases with Bose–Einstein condensates, or unconventional liquids such as helium-4 or helium-3 at sufficiently low temperature. For more than 50 years it was thus unclear whether the supersolid state can exist.
Experiments using helium
While several experiments yielded negative results, in the 1980s, John Goodkind discovered the first anomaly in a solid by using ultrasound. Inspired by his observation, in 2004 Eun-Seong Kim and Moses Chan at Pennsylvania State University saw phenomena which were interpreted as supersolid behavior. Specifically, they observed a non-classical rotational moment of inertia of a torsional oscillator. This observation could not be explained by classical models but was consistent with superfluid-like behavior of a small percentage of the helium atoms contained within the oscillator.
This observation triggered a large number of follow-up studies to reveal the role played by crystal defects or helium-3 impurities. Further experimentation has cast some doubt on the existence of a true supersolid in helium. Most importantly, it was shown that the observed phenomena could be largely explained due to changes in the elastic properties of the helium. In 2012, Chan repeated his original experiments with a new apparatus that was designed to eliminate any such contributions. In this experiment, Chan and his coauthors found no evidence of supersolidity.
Experiments using ultracold quantum gases
In 2017, two research groups from ETH Zurich and from MIT reported on the creation of an ultracold quantum gas with supersolid properties. The Zurich group placed a Bose–Einstein condensate inside two optical resonators, which enhanced the atomic interactions until they started to spontaneously crystallize and form a solid that maintains the inherent superfluidity of Bose–Einstein condensates. This setting realises a special form of a supersolid, the so-called lattice supersolid, where atoms are pinned to the sites of an externally imposed lattice structure. The MIT group exposed a Bose–Einstein condensate in a double-well potential to light beams that created an effective spin–orbit coupling. The interference between the atoms on the two spin–orbit coupled lattice sites gave rise to a characteristic density modulation.
In 2019, three groups from Stuttgart, Florence, and Innsbruck observed supersolid properties in dipolar Bose–Einstein condensates formed from lanthanide atoms. In these systems, supersolidity emerges directly from the atomic interactions, without the need for an external optical lattice. This facilitated also the direct observation of superfluid flow and hence the definitive proof for the existence of the supersolid state of matter.
In 2021, confocal cavity quantum electrodynamics with a Bose–Einstein condensate was used to create a supersolid that possesses a key property of solids, vibration. That is, a supersolid was created that possesses lattice phonons with a Goldstone mode dispersion exhibiting a 16 cm/s speed of sound.
In 2021, dysprosium was used to create a 2-dimensional supersolid quantum gas, in 2022, the same team created a supersolid disk in a round trap and in 2024 they reported the observation of quantum vortices in the supersolid phase
Theory
In most theories of this state, it is supposed that vacancies – empty sites normally occupied by particles in an ideal crystal – lead to supersolidity. These vacancies are caused by zero-point energy, which also causes them to move from site to site as waves. Because vacancies are bosons, if such clouds of vacancies can exist at very low temperatures, then a Bose–Einstein condensation of vacancies could occur at temperatures less than a few tenths of a Kelvin. A coherent flow of vacancies is equivalent to a "superflow" (frictionless flow) of particles in the opposite direction. Despite the presence of the gas of vacancies, the ordered structure of a crystal is maintained, although with less than one particle on each lattice site on average. Alternatively, a supersolid can also emerge from a superfluid. In this situation, which is realised in the experiments with atomic Bose–Einstein condensates, the spatially ordered structure is a modulation on top of the superfluid density distribution.
| Physical sciences | States of matter | Physics |
431637 | https://en.wikipedia.org/wiki/Cockle%20%28bivalve%29 | Cockle (bivalve) | A cockle is an edible marine bivalve mollusc. Although many small edible bivalves are loosely called cockles, true cockles are species in the family Cardiidae.
True cockles live in sandy, sheltered beaches throughout the world. The distinctive rounded shells are bilaterally symmetrical, and are heart-shaped when viewed from the end. Numerous radial, evenly spaced ribs are a feature of the shell in most but not all genera (for an exception, see the genus Laevicardium, the egg cockles, which have very smooth shells).
The shell of a cockle is able to close completely (i.e., there is no "gap" at any point around the edge). Though the shell of a cockle may superficially resemble that of a scallop because of the ribs, cockles can be distinguished from scallops morphologically in that cockle shells lack "auricles" (triangular ear-shaped protrusions near the hinge line) and scallop shells lack a pallial sinus. Behaviorally, cockles live buried in sediment, whereas scallops either are free-living and will swim into the water column to avoid a predator, or in some cases live attached by a byssus to a substrate.
The mantle has three apertures (inhalant, exhalant, and pedal) for siphoning water and for the foot to protrude. Cockles typically burrow using the foot, and feed by filtering plankton from the surrounding water. Cockles are capable of "jumping" by bending and straightening the foot. As is the case in many bivalves, cockles display gonochorism (the sex of an individual varies according to conditions), and some species reach maturity rapidly.
The common name "cockle" is also given by seafood sellers to a number of other small, edible marine bivalves which have a somewhat similar shape and sculpture, but are in other families such as the Veneridae (Venus clams) and the ark clams (Arcidae). Cockles in the family Cardiidae are sometimes referred to as "true cockles" to distinguish them from these other species.
Species
There are more than 205 living species of cockles, with many more fossil forms.
The common cockle (Cerastoderma edule) is widely distributed around the coastlines of Northern Europe, with a range extending west to Ireland, the Barents Sea in the north, Norway in the east, and as far south as Senegal.
The dog cockle, Glycymeris glycymeris, has a similar range and habitat to the common cockle, but is not at all closely related, being in the family Glycymerididae. The dog cockle is edible, but due to its toughness when cooked it is generally not eaten, although a process is being developed to solve this problem.
The blood cockle, Tegillarca granosa (not related to the true cockles, instead in the ark clam family, Arcidae), is extensively cultured from southern Korea to Malaysia.
Genera
Living genera within the family Cardiidae include:
Acanthocardia Gray, 1851
Acrosterigma Dall, 1900
Adacna Eichwald, 1838
Afrocardium Tomlin, 1931
Americardia Stewart, 1930
Apiocardia Olsson, 1961
Bucardium Gray, 1853
Cardium Linnaeus, 1758
Cerastoderma Poli, 1795
Ciliatocardium Kafanov, 1974
Clinocardium Keen, 1936
Corculum Röding, 1798
Ctenocardia H. Adams & A. Adams, 1857
Dallocardia Stewart, 1930
Didacna Eichwald, 1838
Dinocardium Dall, 1900
Discors Deshayes, 1858
Europicardium Popov, 1977
Fragum Röding, 1798
Freneixicardia J. A. Schneider, 2002
Frigidocardium Habe, 1951
Fulvia Gray, 1853
Glans Megerle von Mühlfeld, 1811
Goethemia Lambiotte, 1979
Hippopus Lamarck, 1799
Hypanis Pander in Menetries, 1832
Keenaea Habe, 1951
Keenocardium Kafanov, 1974
Laevicardium Swainson, 1840
Lophocardiium P. Fischer, 1887
Lunulicardia Gray, 1853
Lyrocardium Meek, 1876
Maoricardium Marwick, 1944
Microcardium Keen, 1937
Microfragum Habe, 1951
Monodacna Eichwald, 1838
Nemocardium Meek, 1876
Papillicardium Sacco, 1899
Papyridea Swainson, 1840
Parvicardium Monterosato, 1884
Pratulum Iredale, 1924
Procardium ter Poorten & La Perna, 2017
Pseudofulvia Vidal & Kirkendale, 2007
Ringicardium
Serripes Gould, 1841
Trachycardium Mörch, 1853
Tridacna Bruguière, 1797, the "giant clams"
Trigoniocardium
Vasticardium Iredale, 1927
Vepricardium Iredale, 1929
Extinct genera
† Acobaecardium Paramonova, 1986
† Agnocardia Stewart, 1930
† Aktschagylocardium Danukalova, 1996
† Andrusovicardium Paramonova, 1986
† Anechinocardium Hickman, 2015
† Apscheronia Andrusov, 1903
† Arcicardium P. Fischer, 1887
† Arpadicardium Eberzin, 1947
† Austrocardium Freneix & Grant-Mackie, 1978
† Avicardium V. P. Kolesnikov, 1950
† Avicularium Gray, 1853
† Aviculocardium Bagdasarian, 1978
† Bosphoricardium Eberzin, 1947
† Budmania Brusina, 1897
† Byssocardium Tournouër, 1882
† Caladacna Andrusov, 1917
† Caspicardium Astaf'yeva, 1955
† Chartoconcha Andrusov, 1907
† Chokrakia S. V. Popov, 2001
† Dacicardium Papaianopol, 1975
† Didacnoides Astaf'yeva, 1960
† Didacnomya Andrusov, 1923 (uncertain, unassessed)
† Digressodacna Davitashvili & Kitovani, 1964
† Diversicostata Vassoevich & Eberzin, 1930
† Ecericardium Eberzin, 1947
† Eoprosodacna Davitashvili, 1934
† Ethmocardium White, 1880
† Euxinicardium Eberzin, 1947
† Gilletella Marinescu, 1973
† Goniocardium Vasseur, 1880
† Granocardium Gabb, 1869
† Habecardium Glibert & van de Poel, 1970
† Hedecardium Marwick, 1944
† Hellenicardium S. V. Popov & Nevesskaja, 2000
† Horiodacna Stefanescu, 1896
† Integricardium Rollier, 1912
† Korobkoviella Merklin, 1974
† Kubanocardium Muskhelishvili, 1965
† Lahillia Cossmann, 1899
† Limnodacna Eberzin, 1936
† Limnopagetia Schlickum, 1963
† Limnopappia Schlickum, 1962
† Loxocardium Cossmann, 1886
† Luxuridacna Papaianopol, 1980
† Lymnocardium Stoliczka, 1870
† Merklinicardium S. V. Popov, 1982 (uncertain, unassessed)
† Metadacna Eberzin, 1959
† Miricardium Paramonova, 1986
† Moquicardium Eberzin, 1947
† Myocardia Vest, 1861 (uncertain, unassessed)
† Nargicardium Eberzin, 1947
† Obsoletiformes Kojumdgieva, 1969
† Omanidacna Harzhauser & Mandic, 2008
† Oraphocardium Eberzin, 1949
† Orthocardium Tremlett, 1950
† Oxydacna Davitashvili, 1930
† Pachydacna Eberzin, 1955
† Pannonicardium Stevanović, 1951
† Panticapaea Andrusov, 1923
† Papyrocardium Gabuniya, 1953 (uncertain, unassessed)
† Paradacna Andrusov, 1909
† Parapscheronia Eberzin, 1955
† Parvidacna Stevanović, 1950
† Phyllocardium P. Fischer, 1887
† Plagiocardium Cossmann, 1886
† Plagiodacna Andrusov, 1903
† Plagiodacnopsis Andrusov, 1923
† Planacardium Paramonova, 1971
† Plicatiformes Kojumdgieva, 1969
† Pontalmyra Stefanescu, 1896
† Prionopleura Eberzin, 1949
† Prophyllicardium Jekelius, 1944 (uncertain, unassessed)
† Prosochiasta Eberzin, 1959
† Prosodacna Tournouër, 1882
† Prosodacnomya Eberzin, 1959
† Protocardia Beyrich, 1845
† Protoplagiodacna Stevanović, 1978
† Pseudocatillus Andrusov, 1903
† Pteradacna Andrusov, 1907
† Raricardium Paramonova, 1986
† Replidacna Jekelius, 1944
† Schedocardia Stewart, 1930
† Schirvanicardium Andreescu, 1974
† Stylodacna Stefanescu, 1896
† Submonodacna Livental, 1931
† Tauricardium Eberzin, 1947
† Tschaudia Davitashvili & Kitovani, 1964
† Turcmena G. I. Popov, 1956
† Uniocardium Capellini, 1880 (uncertain, unassessed)
† Yokoyamaina Hayami, 1958
† Zamphiridacna Motaş, 1974
Gallery
In cuisine and culture
Cockles are a popular type of edible shellfish in both Eastern and Western cooking. They are collected by raking them from the sands at low tide. However, collecting cockles is hard work and, as seen from the Morecambe Bay disaster, in which 23 people died, can be dangerous if local tidal conditions are not carefully watched.
In England and Wales, , people are permitted to collect 5 kg of cockles for personal use. Those wishing to collect more than this are deemed to be engaging in commercial fishing and are required to obtain a permit from the Inshore Fisheries and Conservation Authority.
Cockles are a street food in Cambodia where it is usually steamed or boiled and served with a dipping sauce consisting of crushed peppercorns, salt and lime juice.
Cockles are sold freshly cooked as a snack in the United Kingdom, particularly in those parts of the British coastline where cockles are abundant. Boiled, then seasoned with malt vinegar and white pepper, they can be bought from seafood stalls, which also often have for sale mussels, whelks, jellied eels, crabs and shrimp. Cockles are also available pickled in jars, and more recently, have been sold in sealed packets (with vinegar) containing a plastic two-pronged fork. A meal of cockles fried with bacon, served with laverbread, is known as a traditional Welsh breakfast.
Boiled cockles (sometimes grilled) are sold at many hawker centres in Southeast Asia, and are used in laksa, char kway teow and steamboat. They are called kerang in Malay and see hum in Cantonese.
In Japan, the Japanese egg cockle (Laevicardium laevigatum) is used to create torigai sushi.
A study conducted in England in the early 1980s showed a correlation between the consumption of cockles, presumed to be incorrectly processed, and an elevated local occurrence of hepatitis.
Cockles are an effective bait for a wide variety of sea fishes. The folk song "Molly Malone" is also known as "Cockles and Mussels" because the title character's sale of the two foods is referred to in the song's refrain. The shells of cockles are mentioned in the English nursery rhyme "Mary, Mary, Quite Contrary". Cockles are also eaten by the indigenous peoples of North America.
Alternative meanings
The common English phrase "it warms the cockles of my heart", is used to mean that a feeling of deep-seated contentment has been generated.
Differing derivations of this phrase have been proposed, either directly from the perceived heart-shape of a cockleshell, or indirectly (the scientific name for the type genus of the family is Cardium, from the Latin for heart), or from the Latin diminutive of the word heart, corculum. Another proposed derivation is from the Latin for the ventricles of the heart, cochleae cordis, where the second word is an inflected form of cor, heart, while cochlea is the Latin for snail.
| Biology and health sciences | Bivalvia | Animals |
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