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725585 | https://en.wikipedia.org/wiki/10%20Hygiea | 10 Hygiea | 10 Hygiea is a major asteroid located in the main asteroid belt. With a mean diameter of between km and a mass estimated to be 3% of the total mass of the belt, it is the fourth-largest asteroid in the Solar System by both volume and mass, and is the largest of the C-type asteroids (dark asteroids with a carbonaceous surface) in classifications that use G type for 1 Ceres. It is very close to spherical, apparently because it had re-accreted after the disruptive impact that produced the large Hygiean family of asteroids.
Observation
Despite its size, Hygiea appears very dim when observed from Earth. This is due to its dark surface and its position in the outer main belt. For this reason, six smaller asteroids were observed before Annibale de Gasparis discovered Hygiea on 12 April 1849. At most oppositions, Hygiea has a magnitude that is four magnitudes dimmer than Vesta's, and observing it typically requires at least a telescope. However, while at a perihelic opposition, it can be observed just with 10×50 binoculars as Hygiea would have a magnitude of +9.1.
Discovery and name
On 12 April 1849, in Naples, Italy, astronomer Annibale de Gasparis (age 29) discovered Hygiea. It was the first of his nine asteroid discoveries. The director of the Naples observatory, Ernesto Capocci, named the asteroid. He chose to call it Igea Borbonica ("Bourbon Hygieia"), after the Greek goddess of health, daughter of Asclepius, and in honor of the ruling family of the Kingdom of the Two Sicilies where Naples was located.
In 1852, John Russell Hind wrote that "it is universally termed Hygiea, the unnecessary appendage 'Borbonica' being dropped." The English form is an irregular spelling of Greek Hygieia or Hygeia (Latin Hygea or Hygia).
Symbol
The intended astronomical symbol for Hygiea was a zeta-shaped serpent crowned with a star, in the pipeline for Unicode 17.0 as U+1F779 the serpent and serpent drinking from a bowl are traditional symbols of the goddess Hygieia (cf. U+1F54F 🕏). In later years it was substituted with a rod of Asclepius: (a serpent twined around a staff, U+2695 ⚕), confusing Hygieia with her masculine counterpart. These symbols are now both largely obsolete. In this century, 10 Hygiea has seen some minor astrological use, and its symbol was confused once again, with Asclepsius's rod replaced by Mercury's caduceus: , though in a more elaborate form (U+2BDA ⯚) than the symbol of the planet Mercury. The caduceus has long been mistaken for the rod of Asclepius.
Physical characteristics
Observations taken with the Very Large Telescope's SPHERE imager in 2017 and 2018 revealed that Hygiea is nearly spherical and is close to a hydrostatic equilibrium shape.
Based on spectral evidence, Hygiea's surface is thought to consist of primitive carbonaceous materials similar to those found in carbonaceous chondrite meteorites. Aqueous alteration products have been detected on its surface, which could indicate the presence of water ice in the past which was heated sufficiently to melt. The primitive present surface composition indicates that Hygiea had not melted during the early period of Solar System formation. However, observations suggest Hygiea suffered a major collision early in its history that completely disrupted it, with its present spherical shape due to re-accretion of the disrupted material. No deep basins are visible in VLT images, indicating that any large craters that formed after re-accretion must have flat floors, consistent with an icy C-type composition.
In images taken with the VLT in 2017, a bright surface feature is visible, as well as at least two dark craters, which have been informally named Serpens and Calix after the Latin words for 'snake' and 'cup', respectively. Serpens has a diameter of 180 km, Calix of 90 km.
Hygiea is the largest of the class of dark C-type asteroids that are dominant in the outer asteroid belt, beyond the Kirkwood gap at 2.82 AU. Its mean diameter . Hygiea is close to spherical, with an axis ratio of that is consistent with a MacLaurin ellipsoid. Aside from being the smallest of the "big four", Hygiea has a relatively low density of , comparable to Ceres (2.16) and the larger icy satellites of the Solar System (Ganymede 1.94, Callisto 1.83, Titan 1.88, Triton 2.06) rather than to Pallas () or Vesta (3.45).
Although it is the largest body in its region, due to its dark surface and farther-than-average distance from the Sun, Hygiea appears very dim when observed from Earth. In fact, it is the third dimmest of the first twenty-three asteroids discovered, with only 13 Egeria and 17 Thetis having lower mean opposition magnitudes. At most oppositions, Hygiea has a magnitude of around +10.2, which is as much as four orders fainter than Vesta, and observation calls for at least a telescope to resolve. However, at a perihelic opposition, Hygiea can reach +9.1 magnitude and may just be resolvable with 10 × 50 binoculars, unlike the next two largest asteroids in the asteroid belt, 704 Interamnia and 511 Davida, which are always beyond binocular visibility.
A total of 17 stellar occultations by Hygiea have been tracked by Earth-based astronomers, including two (in 2002 and 2014) that were seen by a large number of observers. The observations have been used to constrain Hygiea's size, shape and rotation axis. The Hubble Space Telescope has resolved the asteroid and ruled out the presence of any orbiting companions larger than about in diameter.
Orbit and rotation
Orbiting at an average of 3.14 AU from the Sun, Hygiea is the most distant of the "big four" asteroids. It lies closer to the ecliptic as well, with an orbital inclination of 4°. Its orbit is less circular than those of Ceres or Vesta, with an eccentricity of around 0.12. Its perihelion is at a quite similar longitude to those of Vesta and Ceres, though its ascending and descending nodes are opposite to the corresponding ones for those objects. Although its perihelion is extremely close to the mean distance of Ceres and Pallas, a collision between Hygiea and its larger companions is impossible because at that distance they are always on opposite sides of the ecliptic. In 2056, Hygiea will pass 0.025 AU from Ceres, and then in 2063, Hygiea will pass 0.020 AU from Pallas. At aphelion Hygiea reaches out to the extreme edge of the asteroid belt at the perihelia of the Hilda family, which is in a 3:2 orbital resonance with Jupiter.
As one of the most massive asteroids, Hygiea is used by the Minor Planet Center to calculate perturbations.
Hygiea is in an unstable three-body mean motion resonance with Jupiter and Saturn. The computed Lyapunov time for this asteroid is 30,000 years, indicating that it occupies a chaotic orbit that will change randomly over time because of gravitational perturbations by the planets. It is the lowest numbered asteroid in such a resonance (the next lowest numbered being 70 Panopaea).
Hygiea has a rotation period of 13.83 hours. Its single-peaked light curve has an amplitude of 0.27 mag, which is largely attributed to albedo variations. Hygiea's north pole points towards ecliptic longitude and ecliptic latitude , which gives an axial tilt of 119° with respect to the ecliptic.
Hygiea family
Hygiea is the main member of the Hygiean asteroid family that constitutes about 1% of asteroids in the main belt. The family was formed when an object with a diameter of about 100 km collided with proto-Hygiea about 2 billion years ago. Because the impact craters on Hygiea today are too small to contain the volume of ejected material, it is thought that Hygiea was completely disrupted by the impact and that the majority of the debris recoalesced after the pieces that formed the rest of the family had escaped. Hygiea contains almost all the mass (over 98%) of the family.
| Physical sciences | Solar System | Astronomy |
725992 | https://en.wikipedia.org/wiki/Blinking | Blinking | Blinking is a bodily function; it is a semi-autonomic rapid closing of the eyelid. A single blink is determined by the forceful closing of the eyelid or inactivation of the levator palpebrae superioris and the activation of the palpebral portion of the orbicularis oculi, not the full open and close. It is an essential function of the eye that helps spread tears across and remove irritants from the surface of the cornea and conjunctiva.
Blinking may have other functions since it occurs more often than necessary just to keep the eye lubricated. Researchers think blinking may help with disengagement of attention; following blink onset, cortical activity decreases in the dorsal network and increases in the default-mode network, associated with internal processing. Blink speed can be affected by elements such as fatigue, eye injury, medication, and disease. The blinking rate is determined by the "blinking center", but it can also be affected by external stimulus.
Some animals, such as tortoises and hamsters, blink their eyes independently of each other. Humans use winking, the blinking of only one eye, as a form of body language.
Function and anatomy
Blinking provides moisture to the eye by irrigation using tears and a lubricant the eyes secrete. The eyelid provides suction across the eye from the tear duct to the entire eyeball to keep it from drying out.
Blinking also protects the eye from irritants. Eyelashes are hairs which grow from the edges of the upper and lower eyelids that create a line of defense against dust and other elements to the eye. The eyelashes catch most of these irritants before they reach the eyeball.
There are multiple muscles that control reflexes of blinking. The main muscles, in the upper eyelid, that control the opening and closing are the orbicularis oculi and levator palpebrae superioris muscle. The orbicularis oculi closes the eye, while the contraction of the levator palpebrae muscle opens the eye. The Müller's muscle, or the superior tarsal muscle, in the upper eyelid and the inferior palpebral muscle in the lower 3 eyelid are responsible for widening the eyes. These muscles are not only imperative in blinking, but they are also important in many other functions such as squinting and winking. The inferior palpebral muscle is coordinated with the inferior rectus to pull down the lower lid when one looks down.
The correlation between human eyelid blink behavior and psychological stress was also demonstrated by means of a laboratory study. Lying may affect the rate of blinking.
Blinking is used for communication in humans, some primates, in human interactions with cats, and by female concave-eared torrent frogs to initiate mating with males.
Central nervous system's control
Though one may think that the stimulus triggering blinking is dry or irritated eyes, it is most likely that it is controlled by a "blinking center" of the globus pallidus of the lenticular nucleus—a body of nerve cells between the base and outer surface of the brain. Nevertheless, external stimuli can contribute. The orbicularis oculi is a facial muscle; therefore its actions are translated by the facial nerve root. The levator palpebrae superioris' action is sent through the oculomotor nerve. The duration of a blink is
on average 100–150 milliseconds according to UCL researcher and between 100 and 400 ms according to the Harvard Database of Useful Biological Numbers. Closures in excess of 1000 ms were defined as microsleeps.
Greater activation of dopaminergic pathways dopamine production in the striatum is associated with a higher rate of spontaneous eye blinking. Conditions in which there is reduced dopamine availability such as Parkinson's disease have reduced eye blink rate, while conditions in which it is raised such as schizophrenia have an increased rate. Blink rate is associated with dopamine-related executive function and creativity.
Evolutionary origins
Blinking is present in all major tetrapod crown groups. The soft tissues involved in blinking have not been preserved in the fossil record, but study of mudskippers (a group of amphibious fish species that evolved blinking independently from other tetrapod species, but for similar purposes), suggest that blinking (which involves the eye retracting in mudskippers) may have arose in response to selective pressures upon species shifting from aquatic to terrestrial habitats. For example, compared to an aquatic environment, in a terrestrial environment, the corneal cells must be kept moist such that vital substances like oxygen can more easily diffuse into them, detritus may adhere to the eye in dry conditions, and objects may move towards the eye at faster and more dangerous speeds in air than in water. Additionally, when at their fully aquatic juvenile stage of development, their eyes are not in the positioning with which they blink, but as adults, their eyes elevate to a position that can blink, which they do when they are not submerged or bump into a surface, suggesting blinking emerged as an adaptation to terrestrial life as opposed to aquatic life.
Early tetrapods in the transition to land, which would later yield all non-mudskipper blinking species, possessed similar characteristics regarding eye positioning that suggest blinking arose in response to aerial vision and terrestrial lifestyle.
Types of blinking
There are three types of blinking.
Spontaneous blink
Spontaneous blinking is done without external stimuli and internal effort. This type of blinking is conducted in the pre-motor brain stem and happens without conscious effort, like breathing and digestion.
Reflex blink
A reflex blink occurs in response to an external stimulus, such as contact with the cornea or objects that appear rapidly in front of the eye. A reflex blink is not necessarily a conscious blink either; however it does happen faster than a spontaneous blink. Reflex blink may occur in response to tactile stimuli (e.g., corneal, eyelash, skin of eyelid, contact with eyebrow), optical stimuli (e.g. dazzle reflex, or menace reflex) or auditory stimuli (e.g., menace reflex).
Voluntary blink
A voluntary blink is a conscious blink, with the use of all 3 divisions of the orbicularis oculi muscle.
Blinking in everyday life
Children
Infants do not blink at the same rate of adults; in fact, infants only blink at an average rate of one or two times in a minute. The reason for this difference is unknown, but it is suggested that infants do not require the same amount of eye lubrication that adults do because their eyelid opening is smaller in relation to adults. Additionally, infants do not produce tears during their first month of life. Infants also get a significant amount more sleep than adults do and, as discussed earlier, fatigued eyes blink more. However, throughout childhood the blink rate increases, and by adolescence, it is usually equivalent to that of adults.
Adults
There have been mixed results when studying gender-dependent differences in blinking rates, with results varying from the women's rate nearly doubling the men's to no significant difference between them. In addition, women using oral contraceptives blink 32% more often than other women on average for unknown reasons. Generally, between each blink is an interval of 2–10 seconds; actual rates vary by individual, averaging around 17 blinks per minute in a laboratory setting. However, when the eyes are focused on an object for an extended period of time, such as when reading, the rate of blinking decreases to about 4 to 5 times per minute. This is the major reason that eyes dry out and become fatigued when reading.
When the eyes dry out or become fatigued due to reading on a computer screen, it can be an indication of computer vision syndrome. Computer vision syndrome can be prevented by taking regular breaks, focusing on objects far from the screen, having a well-lit workplace, or using a blink reminder application. Studies suggest that adults can learn to maintain a healthy blinking rate while reading or looking at a computer screen using biofeedback.
Eye blinking can be a criterion for diagnosing medical conditions. For example, excessive blinking may help to indicate the onset of Tourette syndrome, strokes or disorders of the nervous system. A reduced rate of blinking is associated with Parkinson's disease.
| Biology and health sciences | Basics | Biology |
726049 | https://en.wikipedia.org/wiki/Pharmacodynamics | Pharmacodynamics | Pharmacodynamics (PD) is the study of the biochemical and physiologic effects of drugs (especially pharmaceutical drugs). The effects can include those manifested within animals (including humans), microorganisms, or combinations of organisms (for example, infection).
Pharmacodynamics and pharmacokinetics are the main branches of pharmacology, being itself a topic of biology interested in the study of the interactions of both endogenous and exogenous chemical substances with living organisms.
In particular, pharmacodynamics is the study of how a drug affects an organism, whereas pharmacokinetics is the study of how the organism affects the drug. Both together influence dosing, benefit, and adverse effects. Pharmacodynamics is sometimes abbreviated as PD and pharmacokinetics as PK, especially in combined reference (for example, when speaking of PK/PD models).
Pharmacodynamics places particular emphasis on dose–response relationships, that is, the relationships between drug concentration and effect. One dominant example is drug-receptor interactions as modeled by
L + R <=> LR
where L, R, and LR represent ligand (drug), receptor, and ligand-receptor complex concentrations, respectively. This equation represents a simplified model of reaction dynamics that can be studied mathematically through tools such as free energy maps.
Basics
There are four principal protein targets with which drugs can interact:
Enzymes (e.g. neostigmine and acetyl cholinesterase)
Inhibitors
Inducers
Activators
Membrane carriers [Reuptake vs Efflux] (e.g. tricyclic antidepressants and catecholamine uptake-1)
Enhancer (RE)
Inhibitor (RI)
Releaser (RA)
Ion channels (e.g. nimodipine and voltage-gated Ca2+ channels)
Blocker
Opener
Receptor (e.g. Listed in table below)
Agonists can be full, partial or inverse.
Antagonists can be competitive, non-competitive, or uncompetive.
Allosteric modulator can have 3 effects within a receptor. One is its capability or incapability to activate a receptor (2 possibilities). The other two are agonist affinity and efficacy. They may be increased, decreased or unaffected (3 and 3 possibilities).
NMBD = neuromuscular blocking drugs; NMDA = N-methyl-d-aspartate; EGF = epidermal growth factor.
Effects on the body
The majority of drugs either
There are 7 main drug actions:
stimulating action through direct receptor agonism and downstream effects
depressing action through direct receptor agonism and downstream effects (ex.: inverse agonist)
blocking/antagonizing action (as with silent antagonists), the drug binds the receptor but does not activate it
stabilizing action, the drug seems to act neither as a stimulant or as a depressant (ex.: some drugs possess receptor activity that allows them to stabilize general receptor activation, like buprenorphine in opioid dependent individuals or aripiprazole in schizophrenia, all depending on the dose and the recipient)
exchanging/replacing substances or accumulating them to form a reserve (ex.: glycogen storage)
direct beneficial chemical reaction as in free radical scavenging
direct harmful chemical reaction which might result in damage or destruction of the cells, through induced toxic or lethal damage (cytotoxicity or irritation)
Desired activity
The desired activity of a drug is mainly due to successful targeting of one of the following:
Cellular membrane disruption
Chemical reaction with downstream effects
Interaction with enzyme proteins
Interaction with structural proteins
Interaction with carrier proteins
Interaction with ion channels
Ligand binding to receptors:
Hormone receptors
Neuromodulator receptors
Neurotransmitter receptors
General anesthetics were once thought to work by disordering the neural membranes, thereby altering the Na+ influx. Antacids and chelating agents combine chemically in the body. Enzyme-substrate binding is a way to alter the production or metabolism of key endogenous chemicals, for example aspirin irreversibly inhibits the enzyme prostaglandin synthetase (cyclooxygenase) thereby preventing inflammatory response. Colchicine, a drug for gout, interferes with the function of the structural protein tubulin, while digitalis, a drug still used in heart failure, inhibits the activity of the carrier molecule, Na-K-ATPase pump. The widest class of drugs act as ligands that bind to receptors that determine cellular effects. Upon drug binding, receptors can elicit their normal action (agonist), blocked action (antagonist), or even action opposite to normal (inverse agonist).
In principle, a pharmacologist would aim for a target plasma concentration of the drug for a desired level of response. In reality, there are many factors affecting this goal. Pharmacokinetic factors determine peak concentrations, and concentrations cannot be maintained with absolute consistency because of metabolic breakdown and excretory clearance. Genetic factors may exist which would alter metabolism or drug action itself, and a patient's immediate status may also affect indicated dosage.
Undesirable effects
Undesirable effects of a drug include:
Increased probability of cell mutation (carcinogenic activity)
A multitude of simultaneous assorted actions which may be deleterious
Interaction (additive, multiplicative, or metabolic)
Induced physiological damage, or abnormal chronic conditions
Therapeutic window
The therapeutic window is the amount of a medication between the amount that gives an effect (effective dose) and the amount that gives more adverse effects than desired effects. For instance, medication with a small pharmaceutical window must be administered with care and control, e.g. by frequently measuring blood concentration of the drug, since it easily loses effects or gives adverse effects.
Duration of action
The duration of action of a drug is the length of time that particular drug is effective. Duration of action is a function of several parameters including plasma half-life, the time to equilibrate between plasma and target compartments, and the off rate of the drug from its biological target.
Recreational drug use
In recreational psychoactive drug spaces, duration refers to the length of time over which the subjective effects of a psychoactive substance manifest themselves.
Duration can be broken down into 6 parts: (1) total duration (2) onset (3) come up (4) peak (5) offset and (6) after effects. Depending upon the substance consumed, each of these occurs in a separate and continuous fashion.
Total
The total duration of a substance can be defined as the amount of time it takes for the effects of a substance to completely wear off into sobriety, starting from the moment the substance is first administered.
Onset
The onset phase can be defined as the period until the very first changes in perception (i.e. "first alerts") are able to be detected.
Come up
The "come up" phase can be defined as the period between the first noticeable changes in perception and the point of highest subjective intensity. This is colloquially known as "coming up."
Peak
The peak phase can be defined as period of time in which the intensity of the substance's effects are at its height.
Offset
The offset phase can be defined as the amount of time in between the conclusion of the peak and shifting into a sober state. This is colloquially referred to as "coming down."
After effects
The after effects can be defined as any residual effects which may remain after the experience has reached its conclusion. After effects depend on the substance and usage. This is colloquially known as a "hangover" for negative after effects of substances, such as alcohol, cocaine, and MDMA or an "afterglow" for describing a typically positive, pleasant effect, typically found in substances such as cannabis, LSD in low to high doses, and ketamine.
Receptor binding and effect
The binding of ligands (drug) to receptors is governed by the law of mass action which relates the large-scale status to the rate of numerous molecular processes. The rates of formation and un-formation can be used to determine the equilibrium concentration of bound receptors. The equilibrium dissociation constant is defined by:
L + R <=> LR
where L=ligand, R=receptor, square brackets [] denote concentration. The fraction of bound receptors is
Where is the fraction of receptor bound by the ligand.
This expression is one way to consider the effect of a drug, in which the response is related to the fraction of bound receptors (see: Hill equation). The fraction of bound receptors is known as occupancy. The relationship between occupancy and pharmacological response is usually non-linear. This explains the so-called receptor reserve phenomenon i.e. the concentration producing 50% occupancy is typically higher than the concentration producing 50% of maximum response. More precisely, receptor reserve refers to a phenomenon whereby stimulation of only a fraction of the whole receptor population apparently elicits the maximal effect achievable in a particular tissue.
The simplest interpretation of receptor reserve is that it is a model that states there are excess receptors on the cell surface than what is necessary for full effect. Taking a more sophisticated approach, receptor reserve is an integrative measure of the response-inducing capacity of an agonist (in some receptor models it is termed intrinsic efficacy or intrinsic activity) and of the signal amplification capacity of the corresponding receptor (and its downstream signaling pathways). Thus, the existence (and magnitude) of receptor reserve depends on the agonist (efficacy), tissue (signal amplification ability) and measured effect (pathways activated to cause signal amplification). As receptor reserve is very sensitive to agonist's intrinsic efficacy, it is usually defined only for full (high-efficacy) agonists.
Often the response is determined as a function of log[L] to consider many orders of magnitude of concentration. However, there is no biological or physical theory that relates effects to the log of concentration. It is just convenient for graphing purposes. It is useful to note that 50% of the receptors are bound when [L]=Kd .
The graph shown represents the conc-response for two hypothetical receptor agonists, plotted in a semi-log fashion. The curve toward the left represents a higher potency (potency arrow does not indicate direction of increase) since lower concentrations are needed for a given response. The effect increases as a function of concentration.
Multicellular pharmacodynamics
The concept of pharmacodynamics has been expanded to include Multicellular Pharmacodynamics (MCPD). MCPD is the study of the static and dynamic properties and relationships between a set of drugs and a dynamic and diverse multicellular four-dimensional organization. It is the study of the workings of a drug on a minimal multicellular system (mMCS), both in vivo and in silico. Networked Multicellular Pharmacodynamics (Net-MCPD) further extends the concept of MCPD to model regulatory genomic networks together with signal transduction pathways, as part of a complex of interacting components in the cell.
Toxicodynamics
Toxicodynamics (TD) and pharmacodynamics (PD) link a therapeutic agent or toxicant, or toxin (xenobiotic)'s dosage to the features, amount, and time course of its biological action. The mechanism of action is a crucial factor in determining effect and toxicity of the drug, taking in consideration the pharmacokinetic (PK) factors. The sort and extent of altered cellular physiology will depend on the combination of the drug's presence (as established by pharmacokinetic (PK) studies) and/or its mechanism and duration of action (PD). Types of xenobiotic-target interaction can be described either by reversible, irreversible, noncompetitive, and allosteric interaction or agonist, partial agonist, antagonist, and inverse interactions. In vitro, ex vivo, or in vivo studies can be used to assess PD and TD from the molecule to the level of the entire organism.
The mechanism of drug action and adverse drug reaction is either physiochemical property based and biochemical based. Adverse drugs reactions can be classified as either idiosyncratic (type B) or intrinsic (type A). Idiosyncratic toxicity is not dosage dependent and defy the mass-action relationship. Immune-mediated processes are frequently cited as the source of type B reactions. These cannot be accurately described in preclinical research or clinical trials due to their low incidence frequency. Type A reactions are dosage (concentration) dependent. Usually, this kind of side effect is an extension of an ongoing treatment.
Pharmacokinetics and pharmacodynamics are termed toxicokinetics and toxicodynamics in the field of ecotoxicology. Here, the focus is on toxic effects on a wide range of organisms. The corresponding models are called toxicokinetic-toxicodynamic models.
| Biology and health sciences | Fields of medicine | Health |
726748 | https://en.wikipedia.org/wiki/Black-body%20radiation | Black-body radiation | Black-body radiation is the thermal electromagnetic radiation within, or surrounding, a body in thermodynamic equilibrium with its environment, emitted by a black body (an idealized opaque, non-reflective body). It has a specific, continuous spectrum of wavelengths, inversely related to intensity, that depend only on the body's temperature, which is assumed, for the sake of calculations and theory, to be uniform and constant.
A perfectly insulated enclosure which is in thermal equilibrium internally contains blackbody radiation, and will emit it through a hole made in its wall, provided the hole is small enough to have a negligible effect upon the equilibrium. The thermal radiation spontaneously emitted by many ordinary objects can be approximated as blackbody radiation.
Of particular importance, although planets and stars (including the Earth and Sun) are neither in thermal equilibrium with their surroundings nor perfect black bodies, blackbody radiation is still a good first approximation for the energy they emit.
The term black body was introduced by Gustav Kirchhoff in 1860. Blackbody radiation is also called thermal radiation, cavity radiation, complete radiation or temperature radiation.
Theory
Spectrum
Black-body radiation has a characteristic, continuous frequency spectrum that depends only on the body's temperature, called the Planck spectrum or Planck's law. The spectrum is peaked at a characteristic frequency that shifts to higher frequencies with increasing temperature, and at room temperature most of the emission is in the infrared region of the electromagnetic spectrum. As the temperature increases past about 500 degrees Celsius, black bodies start to emit significant amounts of visible light. Viewed in the dark by the human eye, the first faint glow appears as a "ghostly" grey (the visible light is actually red, but low intensity light activates only the eye's grey-level sensors). With rising temperature, the glow becomes visible even when there is some background surrounding light: first as a dull red, then yellow, and eventually a "dazzling bluish-white" as the temperature rises. When the body appears white, it is emitting a substantial fraction of its energy as ultraviolet radiation. The Sun, with an effective temperature of approximately 5800 K, is an approximate black body with an emission spectrum peaked in the central, yellow-green part of the visible spectrum, but with significant power in the ultraviolet as well.
Blackbody radiation provides insight into the thermodynamic equilibrium state of cavity radiation.
Black body
All normal (baryonic) matter emits electromagnetic radiation when it has a temperature above absolute zero. The radiation represents a conversion of a body's internal energy into electromagnetic energy, and is therefore called thermal radiation. It is a spontaneous process of radiative distribution of entropy.
Conversely, all normal matter absorbs electromagnetic radiation to some degree. An object that absorbs all radiation falling on it, at all wavelengths, is called a black body. When a black body is at a uniform temperature, its emission has a characteristic frequency distribution that depends on the temperature. Its emission is called blackbody radiation.
The concept of the black body is an idealization, as perfect black bodies do not exist in nature. However, graphite and lamp black, with emissivities greater than 0.95, are good approximations to a black material. Experimentally, blackbody radiation may be established best as the ultimately stable steady state equilibrium radiation in a cavity in a rigid body, at a uniform temperature, that is entirely opaque and is only partly reflective. A closed box with walls of graphite at a constant temperature with a small hole on one side produces a good approximation to ideal blackbody radiation emanating from the opening.
Blackbody radiation has the unique absolutely stable distribution of radiative intensity that can persist in thermodynamic equilibrium in a cavity. In equilibrium, for each frequency, the intensity of radiation which is emitted and reflected from a body relative to other frequencies (that is, the net amount of radiation leaving its surface, called the spectral radiance) is determined solely by the equilibrium temperature and does not depend upon the shape, material or structure of the body. For a black body (a perfect absorber) there is no reflected radiation, and so the spectral radiance is entirely due to emission. In addition, a black body is a diffuse emitter (its emission is independent of direction).
Blackbody radiation becomes a visible glow of light if the temperature of the object is high enough. The Draper point is the temperature at which all solids glow a dim red, about . At , a small opening in the wall of a large uniformly heated opaque-walled cavity (such as an oven), viewed from outside, looks red; at , it looks white. No matter how the oven is constructed, or of what material, as long as it is built so that almost all light entering is absorbed by its walls, it will contain a good approximation to blackbody radiation. The spectrum, and therefore color, of the light that comes out will be a function of the cavity temperature alone. A graph of the spectral radiation intensity plotted versus frequency(or wavelength) is called the blackbody curve. Different curves are obtained by varying the temperature.
When the body is black, the absorption is obvious: the amount of light absorbed is all the light that hits the surface. For a black body much bigger than the wavelength, the light energy absorbed at any wavelength λ per unit time is strictly proportional to the blackbody curve. This means that the blackbody curve is the amount of light energy emitted by a black body, which justifies the name. This is the condition for the applicability of Kirchhoff's law of thermal radiation: the blackbody curve is characteristic of thermal light, which depends only on the temperature of the walls of the cavity, provided that the walls of the cavity are completely opaque and are not very reflective, and that the cavity is in thermodynamic equilibrium. When the black body is small, so that its size is comparable to the wavelength of light, the absorption is modified, because a small object is not an efficient absorber of light of long wavelength, but the principle of strict equality of emission and absorption is always upheld in a condition of thermodynamic equilibrium.
In the laboratory, blackbody radiation is approximated by the radiation from a small hole in a large cavity, a hohlraum, in an entirely opaque body that is only partly reflective, that is maintained at a constant temperature. (This technique leads to the alternative term cavity radiation.) Any light entering the hole would have to reflect off the walls of the cavity multiple times before it escaped, in which process it is nearly certain to be absorbed. Absorption occurs regardless of the wavelength of the radiation entering (as long as it is small compared to the hole). The hole, then, is a close approximation of a theoretical black body and, if the cavity is heated, the spectrum of the hole's radiation (that is, the amount of light emitted from the hole at each wavelength) will be continuous, and will depend only on the temperature and the fact that the walls are opaque and at least partly absorptive, but not on the particular material of which they are built nor on the material in the cavity (compare with emission spectrum).
The radiance or observed intensity is not a function of direction. Therefore, a black body is a perfect Lambertian radiator.
Real objects never behave as full-ideal black bodies, and instead the emitted radiation at a given frequency is a fraction of what the ideal emission would be. The emissivity of a material specifies how well a real body radiates energy as compared with a black body. This emissivity depends on factors such as temperature, emission angle, and wavelength. However, it is typical in engineering to assume that a surface's spectral emissivity and absorptivity do not depend on wavelength so that the emissivity is a constant. This is known as the gray body assumption.
With non-black surfaces, the deviations from ideal blackbody behavior are determined by both the surface structure, such as roughness or granularity, and the chemical composition. On a "per wavelength" basis, real objects in states of local thermodynamic equilibrium still follow Kirchhoff's Law: emissivity equals absorptivity, so that an object that does not absorb all incident light will also emit less radiation than an ideal black body; the incomplete absorption can be due to some of the incident light being transmitted through the body or to some of it being reflected at the surface of the body.
In astronomy, objects such as stars are frequently regarded as black bodies, though this is often a poor approximation. An almost perfect blackbody spectrum is exhibited by the cosmic microwave background radiation. Hawking radiation is the hypothetical blackbody radiation emitted by black holes, at a temperature that depends on the mass, charge, and spin of the hole. If this prediction is correct, black holes will very gradually shrink and evaporate over time as they lose mass by the emission of photons and other particles.
A black body radiates energy at all frequencies, but its intensity rapidly tends to zero at high frequencies (short wavelengths). For example, a black body at room temperature () with one square meter of surface area will emit a photon in the visible range (390–750 nm) at an average rate of one photon every 41 seconds, meaning that, for most practical purposes, such a black body does not emit in the visible range.
The study of the laws of black bodies and the failure of classical physics to describe them helped establish the foundations of quantum mechanics.
Further explanation
According to the Classical Theory of Radiation, if each Fourier mode of the equilibrium radiation (in an otherwise empty cavity with perfectly reflective walls) is considered as a degree of freedom capable of exchanging energy, then, according to the equipartition theorem of classical physics, there would be an equal amount of energy in each mode. Since there are an infinite number of modes, this would imply infinite heat capacity, as well as a nonphysical spectrum of emitted radiation that grows without bound with increasing frequency, a problem known as the ultraviolet catastrophe.
In the longer wavelengths this deviation is not so noticeable, as and are very small. In the shorter wavelengths of the ultraviolet range, however, classical theory predicts the energy emitted tends to infinity, hence the ultraviolet catastrophe. The theory even predicted that all bodies would emit most of their energy in the ultraviolet range, clearly contradicted by the experimental data which showed a different peak wavelength at different temperatures (see also Wien's law).
Instead, in the quantum treatment of this problem, the numbers of the energy modes are quantized, attenuating the spectrum at high frequency in agreement with experimental observation and resolving the catastrophe. The modes that had more energy than the thermal energy of the substance itself were not considered, and because of quantization modes having infinitesimally little energy were excluded.
Thus for shorter wavelengths very few modes (having energy more than ) were allowed, supporting the data that the energy emitted is reduced for wavelengths less than the wavelength of the observed peak of emission.
Notice that there are two factors responsible for the shape of the graph, which can be seen as working opposite to one another. Firstly, shorter wavelengths have a larger number of modes associated with them. This accounts for the increase in spectral radiance as one moves from the longest wavelengths towards the peak at relatively shorter wavelengths. Secondly, though, at shorter wavelengths more energy is needed to reach the threshold level to occupy each mode: the more energy needed to excite the mode, the lower the probability that this mode will be occupied. As the wavelength decreases, the probability of exciting the mode becomes exceedingly small, leading to fewer of these modes being occupied: this accounts for the decrease in spectral radiance at very short wavelengths, left of the peak. Combined, they give the characteristic graph.
Calculating the blackbody curve was a major challenge in theoretical physics during the late nineteenth century. The problem was solved in 1901 by Max Planck in the formalism now known as Planck's law of blackbody radiation. By making changes to Wien's radiation law (not to be confused with Wien's displacement law) consistent with thermodynamics and electromagnetism, he found a mathematical expression fitting the experimental data satisfactorily. Planck had to assume that the energy of the oscillators in the cavity was quantized, which is to say that it existed in integer multiples of some quantity. Einstein built on this idea and proposed the quantization of electromagnetic radiation itself in 1905 to explain the photoelectric effect. These theoretical advances eventually resulted in the superseding of classical electromagnetism by quantum electrodynamics. These quanta were called photons and the blackbody cavity was thought of as containing a gas of photons. In addition, it led to the development of quantum probability distributions, called Fermi–Dirac statistics and Bose–Einstein statistics, each applicable to a different class of particles, fermions and bosons.
The wavelength at which the radiation is strongest is given by Wien's displacement law, and the overall power emitted per unit area is given by the Stefan–Boltzmann law. So, as temperature increases, the glow color changes from red to yellow to white to blue. Even as the peak wavelength moves into the ultra-violet, enough radiation continues to be emitted in the blue wavelengths that the body will continue to appear blue. It will never become invisible—indeed, the radiation of visible light increases monotonically with temperature. The Stefan–Boltzmann law also says that the total radiant heat energy emitted from a surface is proportional to the fourth power of its absolute temperature. The law was formulated by Josef Stefan in 1879 and later derived by Ludwig Boltzmann. The formula is given, where E is the radiant heat emitted from a unit of area per unit time, T is the absolute temperature, and is the Stefan–Boltzmann constant.
Equations
Planck's law of blackbody radiation
Planck's law states that
where
For a black body surface, the spectral radiance density (defined per unit of area normal to the propagation) is independent of the angle of emission with respect to the normal. However, this means that, following Lambert's cosine law, is the radiance density per unit area of emitting surface as the surface area involved in generating the radiance is increased by a factor with respect to an area normal to the propagation direction. At oblique angles, the solid angle spans involved do get smaller, resulting in lower aggregate intensities.
The emitted energy flux density or irradiance , is related to the photon flux density through
Wien's displacement law
Wien's displacement law shows how the spectrum of blackbody radiation at any temperature is related to the spectrum at any other temperature. If we know the shape of the spectrum at one temperature, we can calculate the shape at any other temperature. Spectral intensity can be expressed as a function of wavelength or of frequency.
A consequence of Wien's displacement law is that the wavelength at which the intensity per unit wavelength of the radiation produced by a black body has a local maximum or peak, , is a function only of the temperature:
where the constant b, known as Wien's displacement constant, is equal to . is the Lambert W function. So is approximately 2898 μm/T, with the temperature given in kelvins. At a typical room temperature of 293 K (20 °C), the maximum intensity is at .
Planck's law was also stated above as a function of frequency. The intensity maximum for this is given by
In unitless form, the maximum occurs when where The approximate numerical solution is . At a typical room temperature of 293 K (20 °C), the maximum intensity is for .
Stefan–Boltzmann law
By integrating over the frequency the radiance (units: power / [area × solid angle] ) is
by using with and with being the Stefan–Boltzmann constant.
On a side note, at a distance d, the intensity per area of radiating surface is the useful expression
when the receiving surface is perpendicular to the radiation.
By subsequently integrating over the solid angle for all azimuthal angle (0 to ) and polar angle from 0 to , we arrive at the Stefan–Boltzmann law: the power emitted per unit area of the surface of a black body is directly proportional to the fourth power of its absolute temperature:
We used
Applications
Human-body emission
The human body radiates energy as infrared light. The net power radiated is the difference between the power emitted and the power absorbed:
Applying the Stefan–Boltzmann law,
where and are the body surface area and temperature, is the emissivity, and is the ambient temperature.
The total surface area of an adult is about , and the mid- and far-infrared emissivity of skin and most clothing is near unity, as it is for most nonmetallic surfaces. Skin temperature is about 33 °C, but clothing reduces the surface temperature to about 28 °C when the ambient temperature is 20 °C. Hence, the net radiative heat loss is about
The total energy radiated in one day is about 8 MJ, or 2000 kcal (food calories). Basal metabolic rate for a 40-year-old male is about 35 kcal/(m2·h), which is equivalent to 1700 kcal per day, assuming the same 2 m2 area. However, the mean metabolic rate of sedentary adults is about 50% to 70% greater than their basal rate.
There are other important thermal loss mechanisms, including convection and evaporation. Conduction is negligible – the Nusselt number is much greater than unity. Evaporation by perspiration is only required if radiation and convection are insufficient to maintain a steady-state temperature (but evaporation from the lungs occurs regardless). Free-convection rates are comparable, albeit somewhat lower, than radiative rates. Thus, radiation accounts for about two-thirds of thermal energy loss in cool, still air. Given the approximate nature of many of the assumptions, this can only be taken as a crude estimate. Ambient air motion, causing forced convection, or evaporation reduces the relative importance of radiation as a thermal-loss mechanism.
Application of Wien's law to human-body emission results in a peak wavelength of
For this reason, thermal imaging devices for human subjects are most sensitive in the 7–14 micrometer range.
Temperature relation between a planet and its star
The blackbody law may be used to estimate the temperature of a planet orbiting the Sun.
The temperature of a planet depends on several factors:
Incident radiation from its star
Emitted radiation of the planet (for example, Earth's infrared glow)
The albedo effect causing a fraction of light to be reflected by the planet
The greenhouse effect for planets with an atmosphere
Energy generated internally by a planet itself due to radioactive decay, tidal heating, and adiabatic contraction due to cooling.
The analysis only considers the Sun's heat for a planet in a Solar System.
The Stefan–Boltzmann law gives the total power (energy/second) that the Sun emits:
where
The Sun emits that power equally in all directions. Because of this, the planet is hit with only a tiny fraction of it. The power from the Sun that strikes the planet (at the top of the atmosphere) is:
where
Because of its high temperature, the Sun emits to a large extent in the ultraviolet and visible (UV-Vis) frequency range. In this frequency range, the planet reflects a fraction of this energy where is the albedo or reflectance of the planet in the UV-Vis range. In other words, the planet absorbs a fraction of the Sun's light, and reflects the rest. The power absorbed by the planet and its atmosphere is then:
Even though the planet only absorbs as a circular area , it emits in all directions; the spherical surface area being . If the planet were a perfect black body, it would emit according to the Stefan–Boltzmann law
where is the temperature of the planet. This temperature, calculated for the case of the planet acting as a black body by setting , is known as the effective temperature. The actual temperature of the planet will likely be different, depending on its surface and atmospheric properties. Ignoring the atmosphere and greenhouse effect, the planet, since it is at a much lower temperature than the Sun, emits mostly in the infrared (IR) portion of the spectrum. In this frequency range, it emits of the radiation that a black body would emit where is the average emissivity in the IR range. The power emitted by the planet is then:
For a body in radiative exchange equilibrium with its surroundings, the rate at which it emits radiant energy is equal to the rate at which it absorbs it:
Substituting the expressions for solar and planet power in equations 1–6 and simplifying yields the estimated temperature of the planet, ignoring greenhouse effect, :
In other words, given the assumptions made, the temperature of a planet depends only on the surface temperature of the Sun, the radius of the Sun, the distance between the planet and the Sun, the albedo and the IR emissivity of the planet.
Notice that a gray (flat spectrum) ball where comes to the same temperature as a black body no matter how dark or light gray.
Effective temperature of Earth
Substituting the measured values for the Sun and Earth yields:
With the average emissivity set to unity, the effective temperature of the Earth is:
or −18.8 °C.
This is the temperature of the Earth if it radiated as a perfect black body in the infrared, assuming an unchanging albedo and ignoring greenhouse effects (which can raise the surface temperature of a body above what it would be if it were a perfect black body in all spectrums). The Earth in fact radiates not quite as a perfect black body in the infrared which will raise the estimated temperature a few degrees above the effective temperature. If we wish to estimate what the temperature of the Earth would be if it had no atmosphere, then we could take the albedo and emissivity of the Moon as a good estimate. The albedo and emissivity of the Moon are about 0.1054 and 0.95 respectively, yielding an estimated temperature of about 1.36 °C.
Estimates of the Earth's average albedo vary in the range 0.3–0.4, resulting in different estimated effective temperatures. Estimates are often based on the solar constant (total insolation power density) rather than the temperature, size, and distance of the Sun. For example, using 0.4 for albedo, and an insolation of 1400 W m−2, one obtains an effective temperature of about 245 K.
Similarly using albedo 0.3 and solar constant of 1372 W m−2, one obtains an effective temperature of 255 K.
Cosmology
The cosmic microwave background radiation observed today is the most perfect blackbody radiation ever observed in nature, with a temperature of about 2.7 K. It is a "snapshot" of the radiation at the time of decoupling between matter and radiation in the early universe. Prior to this time, most matter in the universe was in the form of an ionized plasma in thermal, though not full thermodynamic, equilibrium with radiation.
According to Kondepudi and Prigogine, at very high temperatures (above 1010 K; such temperatures existed in the very early universe), where the thermal motion separates protons and neutrons in spite of the strong nuclear forces, electron-positron pairs appear and disappear spontaneously and are in thermal equilibrium with electromagnetic radiation. These particles form a part of the black body spectrum, in addition to the electromagnetic radiation.
A black body at room temperature () radiates mostly in the infrared spectrum, which cannot be perceived by the human eye, but can be sensed by some reptiles. As the object increases in temperature to about , the emission spectrum gets stronger and extends into the human visual range, and the object appears dull red. As its temperature increases further, it emits more and more orange, yellow, green, and then blue light (and ultimately beyond violet, ultraviolet).
Light bulb
Tungsten filament lights have a continuous black body spectrum with a cooler colour temperature, around , which also emits considerable energy in the infrared range. Modern-day fluorescent and LED lights, which are more efficient, do not have a continuous black body emission spectrum, rather emitting directly, or using combinations of phosphors that emit multiple narrow spectrums.
History
In query 6 of Isaac Newton's Opticks, he states that "Do not black Bodies conceive heat more easily from Light than those of other Colours do, by reason that the Light falling on them is not reflected outwards, but enters into the Bodies, and is often reflected and refracted within them, until it be stifled and lost?", thereby introducing the notion of a black body. In his first memoir, Augustin-Jean Fresnel (1788–1827) responded to a view he extracted from a French translation of Newton's Opticks. He says that Newton imagined particles of light traversing space uninhibited by the caloric medium filling it, and refutes this view (never actually held by Newton) by saying that a black body under illumination would increase indefinitely in heat.
Balfour Stewart
In 1858, Balfour Stewart described his experiments on the thermal radiative emissive and absorptive powers of polished plates of various substances, compared with the powers of lamp-black surfaces, at the same temperature. Stewart chose lamp-black surfaces as his reference because of various previous experimental findings, especially those of Pierre Prevost and of John Leslie. He wrote, "Lamp-black, which absorbs all the rays that fall upon it, and therefore possesses the greatest possible absorbing power, will possess also the greatest possible radiating power." Stewart's statement assumed a general principle: that there exists a body or surface that has the greatest possible absorbing and radiative power for every wavelength and equilibrium temperature.
Stewart was concerned with selective thermal radiation, which he investigated using plates which selectively radiated and absorbed different wavelengths. He discussed the experiments in terms of rays which could be reflected and refracted, and which obeyed the Stokes-Helmholtz reciprocity principle. His research did not consider that properties of rays are dependent on wavelength, and he did not use tools such as prisms or diffraction gratings. His work was quantitative within these constraints. He made his measurements in a room temperature environment, and quickly so as to catch his bodies in a condition near the thermal equilibrium in which they had been prepared.
Gustav Kirchhoff
In 1859, Gustav Robert Kirchhoff reported the coincidence of the wavelengths of spectrally resolved lines of absorption and emission of visible light. Importantly for thermal physics, he also observed that bright lines or dark lines were apparent depending on the temperature difference between emitter and absorber.
Kirchhoff then went on to consider some bodies that emit and absorb heat radiation, in an opaque enclosure or cavity, in equilibrium at a temperature .
Here is used a notation different from Kirchhoff's. Here, the emitting power denotes a dimensioned quantity, the total radiation emitted by a body labeled by index at temperature . The total absorption ratio of that body is dimensionless, the ratio of absorbed to incident radiation in the cavity at temperature . (In contrast with Balfour Stewart's, Kirchhoff's definition of his absorption ratio did not refer in particular to a lamp-black surface as the source of the incident radiation.) Thus the ratio of emitting power to absorptivity is a dimensioned quantity, with the dimensions of emitting power, because is dimensionless. Also here the wavelength-specific emitting power of the body at temperature is denoted by and the wavelength-specific absorption ratio by . Again, the ratio of emitting power to absorptivity is a dimensioned quantity, with the dimensions of emitting power.
In a second report made in 1859, Kirchhoff announced a new general principle or law for which he offered a theoretical and mathematical proof, though he did not offer quantitative measurements of radiation powers. His theoretical proof was and still is considered by some writers to be invalid. His principle, however, has endured: it was that for heat rays of the same wavelength, in equilibrium at a given temperature, the wavelength-specific ratio of emitting power to absorptivity has one and the same common value for all bodies that emit and absorb at that wavelength. In symbols, the law stated that the wavelength-specific ratio has one and the same value for all bodies. In this report there was no mention of black bodies.
In 1860, still not knowing of Stewart's measurements for selected qualities of radiation, Kirchhoff pointed out that it was long established experimentally that for total heat radiation emitted and absorbed by a body in equilibrium, the dimensioned total radiation ratio has one and the same value common to all bodies. Again without measurements of radiative powers or other new experimental data, Kirchhoff then offered a fresh theoretical proof of his new principle of the universality of the value of the wavelength-specific ratio at thermal equilibrium. His fresh theoretical proof was and still is considered by some writers to be invalid.
But more importantly, it relied on a new theoretical postulate of "perfectly black bodies," which is the reason why one speaks of Kirchhoff's law. Such black bodies showed complete absorption in their infinitely thin most superficial surface. They correspond to Balfour Stewart's reference bodies, with internal radiation, coated with lamp-black. They were not the more realistic perfectly black bodies later considered by Planck. Planck's black bodies radiated and absorbed only by the material in their interiors; their interfaces with contiguous media were only mathematical surfaces, capable neither of absorption nor emission, but only of reflecting and transmitting with refraction.
Kirchhoff's proof considered an arbitrary non-ideal body labeled as well as various perfect black bodies labeled . It required that the bodies be kept in a cavity in thermal equilibrium at temperature . His proof intended to show that the ratio was independent of the nature of the non-ideal body, however partly transparent or partly reflective it was.
His proof first argued that for wavelength and at temperature , at thermal equilibrium, all perfectly black bodies of the same size and shape have the one and the same common value of emissive power , with the dimensions of power. His proof noted that the dimensionless wavelength-specific absorptivity of a perfectly black body is by definition exactly 1. Then for a perfectly black body, the wavelength-specific ratio of emissive power to absorptivity is again just , with the dimensions of power. Kirchhoff considered thermal equilibrium with the arbitrary non-ideal body, and with a perfectly black body of the same size and shape, in place in his cavity in equilibrium at temperature . He argued that the flows of heat radiation must be the same in each case. Thus he argued that at thermal equilibrium the ratio was equal to , which may now be denoted . is a continuous function, dependent only on at fixed temperature , and an increasing function of at fixed wavelength . It vanishes at low temperatures for visible wavelengths, which does not depend on the nature of the arbitrary non-ideal body (Geometrical factors, taken into detailed account by Kirchhoff, have been ignored in the foregoing).
Thus Kirchhoff's law of thermal radiation can be stated: For any material at all, radiating and absorbing in thermodynamic equilibrium at any given temperature , for every wavelength , the ratio of emissive power to absorptivity has one universal value, which is characteristic of a perfect black body, and is an emissive power which we here represent by . (For our notation , Kirchhoff's original notation was simply .)
Kirchhoff announced that the determination of the function was a problem of the highest importance, though he recognized that there would be experimental difficulties to be overcome. He supposed that like other functions that do not depend on the properties of individual bodies, it would be a simple function. Occasionally by historians that function has been called "Kirchhoff's (emission, universal) function," though its precise mathematical form would not be known for another forty years, till it was discovered by Planck in 1900. The theoretical proof for Kirchhoff's universality principle was worked on and debated by various physicists over the same time, and later. Kirchhoff stated later in 1860 that his theoretical proof was better than Balfour Stewart's, and in some respects it was so. Kirchhoff's 1860 paper did not mention the second law of thermodynamics, and of course did not mention the concept of entropy which had not at that time been established. In a more considered account in a book in 1862, Kirchhoff mentioned the connection of his law with Carnot's principle, which is a form of the second law.
According to Helge Kragh, "Quantum theory owes its origin to the study of thermal radiation, in particular to the "blackbody" radiation that Robert Kirchhoff had first defined in 1859–1860."
Doppler effect
The relativistic Doppler effect causes a shift in the frequency f of light originating from a source that is moving in relation to the observer, so that the wave is observed to have frequency f''':
where v is the velocity of the source in the observer's rest frame, θ is the angle between the velocity vector and the observer-source direction measured in the reference frame of the source, and c is the speed of light. This can be simplified for the special cases of objects moving directly towards (θ = π) or away (θ = 0) from the observer, and for speeds much less than c.
Through Planck's law the temperature spectrum of a black body is proportionally related to the frequency of light and one may substitute the temperature (T) for the frequency in this equation.
For the case of a source moving directly towards or away from the observer, this reduces to
Here v > 0 indicates a receding source, and v < 0 indicates an approaching source.
This is an important effect in astronomy, where the velocities of stars and galaxies can reach significant fractions of c. An example is found in the cosmic microwave background radiation, which exhibits a dipole anisotropy from the Earth's motion relative to this blackbody radiation field.
| Physical sciences | Thermodynamics | null |
726915 | https://en.wikipedia.org/wiki/Alkaline%20battery | Alkaline battery | An alkaline battery (IEC code: L) is a type of primary battery where the electrolyte (most commonly potassium hydroxide) has a pH value above 7. Typically these batteries derive energy from the reaction between zinc metal and manganese dioxide.
Compared with zinc–carbon batteries of the Leclanché cell or zinc chloride types, alkaline batteries have a higher energy density and longer shelf life, yet provide the same voltage.
The alkaline battery gets its name because it has an alkaline electrolyte of potassium hydroxide (KOH) instead of the acidic ammonium chloride (NH4Cl) or zinc chloride (ZnCl2) electrolyte of the zinc–carbon batteries. Other battery systems also use alkaline electrolytes, but they use different active materials for the electrodes.
Alkaline batteries account for 80% of manufactured batteries in the US and over 10 billion individual units produced worldwide. In Japan, alkaline batteries account for 46% of all primary battery sales. In Switzerland, alkaline batteries account for 68%, in the UK 60% and in the EU 47% of all battery sales including secondary types. Alkaline batteries contain zinc (Zn) and manganese dioxide (MnO2) (Health codes 1), which is a cumulative neurotoxin and can be toxic in higher concentrations. However, compared to other battery types, the toxicity of alkaline batteries is moderate.
Alkaline batteries are used in many household items such as Portable media players, digital cameras, toys, flashlights, and radios.
History
Batteries with alkaline (rather than acid) electrolyte were first developed by Waldemar Jungner in 1899, and, working independently, Thomas Edison in 1901. The modern alkaline dry battery, using the zinc/manganese dioxide chemistry, was invented by the Canadian engineer Lewis Urry in the 1950s in Canada before he started working for Union Carbide's Eveready Battery division in Cleveland, OH, building on earlier work by Edison. On October 9, 1957, Urry, Karl Kordesch, and P. A. Marsal filed US patent (2,960,558) for the alkaline battery. It was granted in 1960 and was assigned to the Union Carbide Corporation.
When alkaline batteries were introduced in the late 1960s, their zinc electrodes (in common with the then ubiquitous carbon-zinc cells) had a surface film of mercury amalgam. Its purpose was to control electrolytic action on impurities in the zinc; that unwanted electrolytic action would reduce shelf life and promote leakage. When reductions in mercury content were mandated by various legislatures, it became necessary to greatly improve the purity and consistency of the zinc.
Chemistry
In an alkaline battery, the negative electrode is zinc and the positive electrode is manganese dioxide (MnO2). The alkaline electrolyte of potassium hydroxide (KOH) is not consumed during the reaction (it is regenerated), only the zinc and MnO2 are consumed during discharge. The concentration of alkaline electrolyte of potassium hydroxide remains constant, as there are equal amounts of OH− anions consumed and produced in the two half-reactions occurring at the electrodes.
The two half-reactions are:
Anode (oxidation reaction), negatively charged electrode because accepting from the reductant in the cell:
Cathode (reduction reaction), positively charged electrode because giving to the oxidizer in the cell:
The overall reaction (sum of anodic and cathodic reactions) is:
Capacity
The capacity of an alkaline battery is strongly dependent on the load. An AA-sized alkaline battery might have an effective capacity of at low drain, but at a load of , which is common for digital cameras, the capacity could be as little as . The voltage of the battery declines steadily during use, so the total usable capacity depends on the cutoff voltage of the application.
Voltage
The nominal voltage of a fresh alkaline cell as established by manufacturer standards is . The zero-load voltage of a new alkaline battery ranges from , depending on the purity of the manganese dioxide used and the contents of zinc oxide in the electrolyte. The voltage delivered to a load decreases as the current drawn increases and as the cell discharges. A cell is considered fully discharged when the voltage drops to about . Cells connected in series produce a voltage equal to the sum of the voltages of each cell (e.g., three cells generate about 4.5 V when new).
Current
The amount of electrical current an alkaline battery can deliver is roughly proportional to its physical size. This is a result of decreasing internal resistance as the internal surface area of the cell increases. A rule of thumb is that an AA alkaline battery can deliver without any significant heating. Larger cells, such as C and D cells, can deliver more current. Applications requiring currents of several amperes such as powerful portable audio equipment require D-sized cells to handle the increased load.
In comparison, Lithium-ion and Ni-MH batteries can handle with ease on the standard AA size.
Construction
Alkaline batteries are manufactured in standard cylindrical forms interchangeable with zinc–carbon batteries, and in button forms. Several individual cells may be interconnected to form a true "battery", such as the 9-volt PP3-size battery.
A cylindrical cell is contained in a drawn stainless steel can, which is the cathode connection. The positive electrode mixture is a compressed paste of manganese dioxide with carbon powder added for increased conductivity. The paste may be pressed into the can or deposited as pre-molded rings. The hollow center of the cathode is lined with a separator, which prevents contact of the electrode materials and short-circuiting of the cell. The separator is made of a non-woven layer of cellulose or a synthetic polymer. The separator must conduct ions and remain stable in the highly alkaline electrolyte solution.
The negative electrode is composed of a dispersion of zinc powder in a gel containing the potassium hydroxide electrolyte. The zinc powder provides more surface area for chemical reactions to take place, compared to a metal can. This lowers the internal resistance of the cell. To prevent gassing of the cell at the end of its life, more manganese dioxide is used than required to react with all the zinc. Also, a plastic-made gasket is usually added to increase leakage resistance.
The cell is then wrapped in aluminium foil, a plastic film, or rarely, cardboard, which acts as a final layer of leak protection as well as providing a surface on which logos and labels can be printed.
When describing AAA, AA, C, sub-C and D size cells, the negative electrode is connected to the flat end, and the positive terminal is the end with the raised button. This is usually reversed in button cells, with the flat-ended cylindrical can being the positive terminal.
Recharging of alkaline batteries
Some alkaline batteries are designed to be recharged a few times, and are described as rechargeable alkaline batteries. Attempts to recharge standard alkaline batteries may cause rupture, or the leaking of hazardous liquids that corrode the equipment. However, it is reported that standard alkaline batteries can often be recharged a few times (typically not more than ten), albeit with reduced capacity after each charge; chargers are available commercially. The UK consumer organisation Which? reported that it tested two such chargers with Energizer alkaline batteries, finding that battery capacity dropped on average to 10% of its original value, with huge variations, after two cycles (without stating how depleted they were before recharging) after recharging them two times.
In 2017 Gautam G. Yadav published papers reporting that alkaline batteries made by interleaving the interlayers with copper ions could be recharged for over 6,000 cycles due to the theoretical second electron capacity of manganese dioxide. The energy density of these rechargeable batteries with copper intercalated manganese dioxide is reported to be over , the best among the aqueous-based chemistries. It could be capable of energy densities comparable to lithium-ion (at least ) if zinc utilization in the batteries were improved.
Leaks
Alkaline batteries are prone to leaking potassium hydroxide, a caustic agent that can cause respiratory, eye and skin irritation. The risk of this can be reduced by storing batteries in a dry place and at room temperature. Damage from leakage is mitigated by removing batteries when storing devices. Applying reverse current (such as by recharging disposable-grade cells, or by mixing batteries of different types or state of charge in the same device) can increase the risk of leakage.
All batteries gradually self-discharge (whether installed in a device or not) and dead batteries eventually leak. Extremely high temperatures can also cause batteries to rupture and leak (such as in a car during summer) as well as decrease the shelf life of the battery.
The reason for leaks is that as batteries dischargeeither through usage or gradual self-dischargethe chemistry of the cells changes and some hydrogen gas is generated. This out-gassing increases pressure in the battery. Eventually, the excess pressure either ruptures the insulating seals at the end of the battery, or the outer metal canister, or both. In addition, as the battery ages, its steel outer canister may gradually corrode or rust, which can further contribute to containment failure.
Once a leak has formed due to corrosion of the outer steel shell, potassium hydroxide absorbs carbon dioxide from the air to form a feathery crystalline structure of potassium carbonate that grows and spreads out from the battery over time, following along metal electrodes to circuit boards where it commences oxidation of copper tracks and other components, leading to permanent circuitry damage.
The leaking crystalline growths can also emerge from seams around battery covers to form a furry coating outside the device, that corrodes any objects in contact with the leaking device.
Disposal
Since alkaline batteries were made with less mercury beginning in 1996, alkaline batteries are allowed to be disposed of as regular domestic waste in some locations. However, older alkaline batteries with mercury, and the remaining other heavy metals and corrosive chemicals in all batteries (new and old), still present problems for disposal—especially in landfills. There is also the issue of simplifying the disposal of batteries by excluding them all from domestic waste, so that the most toxic batteries are diverted from general waste streams.
Disposal varies by jurisdiction. For example, the state of California considers all batteries as hazardous waste when discarded, and has banned the disposal of batteries in domestic waste. In Europe, battery disposal is controlled by the WEEE Directive and Battery Directive regulations, and as such alkaline batteries must not be thrown in with domestic waste. In the EU, most stores that sell batteries are required by law to accept old batteries for recycling.
Recycling
The use of disposable batteries increases by 5–6% every year. In the past, used batteries ended up at landfill sites, but in 2004, disposal of alkaline batteries at landfill sites was forbidden by an EU regulation. EU member countries are committed to recycling 50% of alkaline batteries by 2016. The need for recycling thus amounts to per year. The share of alkaline batteries is approximately 80% of the whole.
In the US, only one state, California, requires all alkaline batteries to be recycled. Vermont also has a statewide alkaline battery collection program. In other US states, individuals can purchase battery recycling kits used to ship batteries to recyclers. Some stores such as IKEA also collect alkaline batteries for recycling. However, some chain stores that advertise battery recycling (such as Best Buy) accept rechargeable batteries only, and generally do not accept alkaline batteries.
For recycling, the metals from crushed alkaline batteries are mechanically separated, and the waste black mass is treated chemically to separate zinc, manganese dioxide and potassium hydroxide.
| Technology | Energy storage | null |
19278548 | https://en.wikipedia.org/wiki/Atlas%20%28computer%29 | Atlas (computer) | The Atlas was one of the world's first supercomputers, in use from 1962 (when it was claimed to be the most powerful computer in the world) to 1972. Atlas's capacity promoted the saying that when it went offline, half of the United Kingdom's computer capacity was lost. It is notable for being the first machine with virtual memory (at that time referred to as "one-level store") using paging techniques; this approach quickly spread, and is now ubiquitous.
Atlas was a second-generation computer, using discrete germanium transistors. Atlas was created in a joint development effort among the University of Manchester, Ferranti and Plessey. Two other Atlas machines were built: one for BP and the University of London, and one for the Atlas Computer Laboratory at Chilton near Oxford.
A derivative system was built by Ferranti for the University of Cambridge. Called the Titan, or Atlas 2, it had a different memory organisation and ran a time-sharing operating system developed by Cambridge University Computer Laboratory. Two further Atlas 2s were delivered: one to the CAD Centre in Cambridge (later called CADCentre, then AVEVA), and the other to the Atomic Weapons Research Establishment (AWRE), Aldermaston.
The University of Manchester's Atlas was decommissioned in 1971. The final Atlas, the CADCentre machine, was switched off in late 1976. Parts of the Chilton Atlas are preserved by National Museums Scotland in Edinburgh; the main console itself was rediscovered in July 2014 and is at Rutherford Appleton Laboratory in Chilton, near Oxford.
History
Background
Through 1956 there was a growing awareness that the UK was falling behind the US in computer development. In April, B.W. Pollard of Ferranti told a computer conference that "there is in this country a range of medium-speed computers, and the only two machines which are really fast are the Cambridge EDSAC 2 and the Manchester Mark 2, although both are still very slow compared with the fastest American machines." This was followed by similar concerns expressed in May report to the Department of Scientific and Industrial Research Advisory Committee on High Speed Calculating Machines, better known as the Brunt Committee.
Through this period, Tom Kilburn's team at the University of Manchester had been experimenting with transistor-based systems, building two small machines to test various techniques. This was clearly the way forward, and in the autumn of 1956, Kilburn began canvassing possible customers on what features they would want in a new transistor-based machine. Most commercial customers pointed out the need to support a wide variety of peripheral devices, while the Atomic Energy Authority suggested a machine able to perform an instruction every microsecond, or as it would be known today, 1 MIPS of performance. This later request led to the name of the prospective design, MUSE, for microsecond engine.
The need to support many peripherals and the need to run fast are naturally at odds. A program that processes data from a card reader, for instance, will spend the vast majority of its time waiting for the reader to send in the next bit of data. To support these devices while still making efficient use of the central processing unit (CPU), the new system would need to have additional memory to buffer data and have an operating system that could coordinate the flow of data around the system.
Muse becomes Atlas
When the Brunt Committee heard of new and much faster US designs, the Univac LARC and IBM STRETCH, they were able to gain the attention of the National Research Development Corporation (NRDC), responsible for moving technologies from war-era research groups into the market. Over the next eighteen months, they held numerous meetings with prospective customers, engineering teams at Ferranti and EMI, and design teams at Manchester and the Royal Radar Establishment.
In spite of all this effort, by the summer of 1958, there was still no funding available from the NRDC. Kilburn decided to move things along by building a smaller Muse to experiment with various concepts. This was paid for using funding from the Mark 1 Computer Earnings Fund, which collected funds by renting out time on the University's Mark 1. Soon after the project started, in October 1958, Ferranti decided to become involved. In May 1959 they received a grant of £300,000 from the NRDC to build the system, which would be returned from the proceeds of sales. At some point during this process, the machine was renamed Atlas.
The detailed design was completed by the end of 1959, and the construction of the compilers was proceeding. However, the Supervisor operating system was already well behind. This led to David Howarth, newly hired at Ferranti, expanding the operating system team from two to six programmers. In what is described as a Herculean effort, led by the tireless and energetic Howarth (who completed his Ph.D. in physics at age 22), the team eventually delivered a Supervisor consisting of 35,000 lines of assembler language which had support for multiprogramming to solve the problem of peripheral handling.
Installations
The first Atlas was built up at the university throughout 1962. The schedule was further constrained by the planned shutdown of the Ferranti Mercury machine at the end of December. Atlas met this goal, and was officially commissioned on 7 December by John Cockcroft, director of the AEA. This system had only an early version of Supervisor, and the only compiler was for Autocode. It was not until January 1964 that the final version of Supervisor was installed, along with compilers for ALGOL 60 and Fortran.
By the mid-1960s the original machine was in continual use, based on a 20-hour-per-day schedule, during which time as many as 1,000 programs might be run. Time was split between the University and Ferranti, the latter of which charged £500 an hour to its customers. A portion of this was returned to the University Computer Earnings Fund. In 1969, it was estimated that the computer time received by the University would cost £720,000 if it had been leased on the open market. The machine was shut down on 30 November 1971.
Ferranti sold two other Atlas installations, one to a joint consortium of University of London and BP in 1963, and another to the Atomic Energy Research Establishment (Harwell) in December 1964. The AEA machine was later moved to the Atlas Computer Laboratory at Chilton, a few yards outside the boundary fence of Harwell, which placed it on civilian lands and thus made it much easier to access. This installation grew to be the largest Atlas, containing 48 kWords of 48-bit core memory and 32 tape drives. Time was made available to all UK universities. It was shut down in March 1974.
Titan and Atlas 2
In February 1962, Ferranti gave some parts of an Atlas machine to University of Cambridge, and in return, the University would use these to develop a cheaper version of the system. The result was the Titan machine, which became operational in the summer of 1963. Ferranti sold two more of this design under the name Atlas 2, one to the Atomic Weapons Research Establishment (Aldermaston) in 1963, and another to the government-sponsored Computer Aided Design Center in 1966.
Legacy
Atlas had been designed as a response to the US LARC and STRETCH programs. Both ultimately beat Atlas into official use, LARC in 1961, and STRETCH a few months before Atlas. Atlas was much faster than LARC, about four times, and ran slightly slower than STRETCH - Atlas added two floating-point numbers in about 1.59 microseconds, while STRETCH did the same in 1.38 to 1.5 microseconds. Nevertheless, the head of Ferranti's Software Division, Hugh Devonald, said in 1962: "Atlas is in fact claimed to be the world's most powerful computing system. By such a claim it is meant that, if Atlas and any of its rivals were presented simultaneously with similar large sets of representative computing jobs, Atlas should complete its set ahead of all other computers.". No further sales of LARC were attempted, and it is not clear how many STRETCH machines were ultimately produced.
It was not until 1964's arrival of the CDC 6600 that the Atlas was significantly bested. CDC later stated that it was a 1959 description of Muse that gave CDC ideas that significantly accelerated the development of the 6600 and allowed it to be delivered earlier than originally estimated. This led to it winning a contract for the CSIRO in Australia, which had originally been in discussions to buy an Atlas.
Ferranti was having serious financial difficulties in the early 1960s, and decided to sell the computer division to International Computers and Tabulators (ICT) in 1963. ICT decided to focus on the mid-range market with their ICT 1900 series, a flexible range of machines based on the Canadian Ferranti-Packard 6000.
The Atlas was highly regarded by many in the computer industry. Among its admirers was C. Gordon Bell of Digital Equipment Corporation, who later praised it:
In June 2022 an IEEE Milestone was dedicated to the "Atlas Computer and the Invention of Virtual Memory 1957-1962".
Design
Hardware
The machine had many innovative features, but the key operating parameters were as follows (the store size relates to the Manchester installation; the others were larger):
48-bit word size. A word could hold one floating-point number, one instruction, two 24-bit addresses or signed integers, or eight 6-bit characters.
A fast adder that used novel circuitry to minimise carry propagation time.
24-bit (2 million words, 16 million characters) address space that embraced supervisor ('sacred') store, V-store, fixed store and the user store
16K words of core store (equivalent to 96 KB), featuring interleaving of odd/even addresses
8K words of read-only memory (referred to as the fixed store). This contained the supervisor and extracode routines.
96K words of drum store (eqv. to 576 KB), split across four drums but integrated with the core store using virtual memory. The page size was 512 words, i.e. 3072 bytes.
128 high-speed index registers (B-lines) that could be used for address modification in the mostly double-modified instructions. The register address space also included special registers such as the extracode operand address and the exponent of the floating-point accumulator. Three of the 128 registers were program counter registers: 125 was supervisor (interrupt) control, 126 was extracode control, and 127 was user control. Register 0 always held value 0.
Capability for the addition of (for the time) sophisticated new peripherals such as magnetic tape, including direct memory access (DMA) facilities
Peripheral control through V-store addresses (memory-mapped I/O), interrupts and extracode routines, by reading and writing special wired-in store addresses.
An associative memory (content-addressable memory) of page address registers to determine whether the desired virtual memory location was in core store
Instruction pipelining
Atlas did not use a synchronous clocking mechanism — it was an asynchronous processor — so performance measurements were not easy, but as an example:
Fixed-point register add – 1.59 microseconds
Floating-point add, no modification – 1.61 microseconds
Floating-point add, double modify – 2.61 microseconds
Floating-point multiply, double modify – 4.97 microseconds
Extracode
One feature of the Atlas was "Extracode", a technique that allowed complex instructions to be implemented in software. Dedicated hardware expedited entry to and return from the extracode routine and operand access; also, the code of the extracode routines was stored in ROM, which could be accessed faster than the core store.
The uppermost ten bits of a 48-bit Atlas machine instruction were the operation code. If the most significant bit was set to zero, this was an ordinary machine instruction executed directly by the hardware. If the uppermost bit was set to one, this was an Extracode and was implemented as a special kind of subroutine jump to a location in the fixed store (ROM), its address being determined by the other nine bits. About 250 extracodes were implemented, of the 512 possible.
Extracodes were what would be called software interrupts or traps today. They were used to call mathematical procedures which would have been too inefficient to implement in hardware, for example sine, logarithm, and square root. But about half of the codes were designated as Supervisor functions, which invoked operating system procedures. Typical examples would be "Print the specified character on the specified stream" or "Read a block of 512 words from logical tape N". Extracodes were the only means by which a program could communicate with the Supervisor. Other UK machines of the era, such as the Ferranti Orion, had similar mechanisms for calling on the services of their operating systems.
Software
Atlas pioneered many software concepts still in common use today, including the Atlas Supervisor, "considered by many to be the first recognisable modern operating system".
One of the first high-level languages available on Atlas was named Atlas Autocode, which was contemporary to Algol 60 and created specifically to address what Tony Brooker perceived to be some defects in Algol 60. The Atlas did however support Algol 60, as well as Fortran and COBOL, and ABL (Atlas Basic Language, a symbolic input language close to machine language). Being a university computer it was patronised by a large number of the student population, who had access to a protected machine code development environment.
Several of the compilers were written using the Brooker Morris Compiler Compiler (BMCC), considered to be the first of its type.
It also had a programming language called SPG (System Program Generator). At run time an SPG program could compile more program for itself. It could define and use macros. Its variables were in <angle brackets> and it had a text parser, giving SPG program text a resemblance to Backus–Naur form.
From the outset, Atlas was conceived as a supercomputer that would include a comprehensive operating system. The hardware included specific features that facilitated the work of the operating system. For example, the extracode routines and the interrupt routines each had dedicated storage, registers and program counters; a context switch from user mode to extracode mode or executive mode, or from extracode mode to executive mode, was therefore very fast.
| Technology | Early computers | null |
19278686 | https://en.wikipedia.org/wiki/Armoured%20personnel%20carrier | Armoured personnel carrier | An armoured personnel carrier (APC) is a broad type of armoured military vehicle designed to transport personnel and equipment in combat zones. Since World War I, APCs have become a very common piece of military equipment around the world.
According to the definition in the Treaty on Conventional Armed Forces in Europe, an APC is "an armoured combat vehicle which is designed and equipped to transport a combat infantry squad and which, as a rule, is armed with an integral or organic weapon of less than 20 millimetres calibre." Compared to infantry fighting vehicles (IFVs), which are also used to carry infantry into battle, APCs have less armament and are not designed to provide direct fire support in battle. Infantry units that travel in APCs are known as mechanized infantry. Some militaries also make a distinction between infantry units that use APCs and infantry units that use IFVs, with the latter being known as armoured infantry.
History
One of the first armored vehicles to be used in combat was the Spanish Schneider-Brillié, which saw action in Morocco. It was built from the chassis of a Schneider P2-4000 bus and could carry 12 passengers.
The genesis of the armoured personnel carrier was on the Western Front of World War I. In the later stage of the war, Allied tanks could break through enemy trenches, but the infantry following—who were needed to consolidate the territory acquired—still faced small arms and artillery fire. Without infantry support, the tanks were isolated and more easily destroyed. In response, the British experimented with carrying machine-gun crews in the Mark V* tank, but it was found that the conditions inside the tanks rendered the men unfit for combat.
During World War II, half-tracks like the American M3 and German Sd.Kfz. 251 played a role similar to post-war APCs. British Commonwealth forces relied on the full-tracked Universal Carrier. Over the course of the war, APCs evolved from simple armoured cars with transport capacity to purpose-built vehicles. Obsolete armoured vehicles were also repurposed as APCs, such as the various "Kangaroos" converted from M7 Priest self-propelled guns and from Churchill, M3 Stuart and Ram tanks.
During the Cold War, more specialized APCs were developed. The United States introduced a series of them, including successors to the wartime Landing Vehicle Tracked. The most numerous was the M113 armored personnel carrier, of which more than 80,000 were produced. Western nations have since retired most M113s, replacing them with newer APCs, many of these wheeled. A cold war example of a "Kangaroo" is the heavily armoured Israeli Achzarit, converted from captured T-55s tanks, the concept culminating in the Namer.
Meanwhile, the Warsaw Pact developed their own versions of the APC. The Soviet Union termed theirs the Bronetransporter (), better known as the BTR series. It comprised the BTR-40, BTR-152, BTR-60, BTR-70, BTR-80, and the BTR-90, which as a whole were produced in large numbers. Czechoslovakia and Poland together developed the universal amphibious OT-64 SKOT. The BMP series is termed as infantry fighting vehicles, but it has a designed role of carrying troops to the battlefield. The BMP-1, 2, and 3 all possess the ability to transport troops.
Design
By convention, armoured personnel carriers are not intended to take part in direct-fire battle, but are armed for self-defence and armoured to provide protection from shrapnel and small arms fire.
Mobility
An APC is either wheeled or tracked, or occasionally a combination of the two, as in a half-track.
Wheeled vehicles are typically faster on road and less expensive, however have higher ground pressure which decreases mobility offroad and makes them more likely to become stuck in soft terrains such as mud, snow or sand.
Tracked vehicles typically have lower ground pressure and more maneuverability off-road. Due to the limited service life of their treads, and the wear they cause on roads, tracked vehicles are typically transported over long distances by rail or trucks.
Many APCs are amphibious, meaning they are able to traverse bodies of water. To move in water they will often have propellers or water jets, be propelled by their tracks, or driving on the river bed. Preparing the APC to operate amphibiously usually comprises checking the integrity of the hull and folding down a trim vane in front. Water traverse speed varies greatly between vehicles and is much less than ground speed. The maximum swim speed of the M113 is 3.6 mph (5.8 km/h), about 10% its road speed, and the AAVP-7 can swim at 8.2 mph (13.2 km/h).
Protection
Armoured personnel carriers are typically designed to protect against small arms and artillery fire. Some designs have more protection; the Israeli IDF Namer has as much armour as Merkava main battle tank. Armour is usually composed of steel or aluminium. They will also use ballistic glass.
Many APCs are equipped with CBRN protection, which is intended to provide protection from weapons of mass destruction like poison gas and radioactive/nuclear weapons.
Generally APCs will be lighter and less armoured than tanks or IFVs, often being open topped and featuring doors and windows, as seen in the French VAB.
Weaponry
Armoured personnel carriers are designed primarily for transport and are lightly armed. They may be unarmed, or armed with some combination of light, general-purpose, heavy machine guns, or automatic grenade launchers.
In Western nations, APCs are frequently armed with the .50 calibre M2 Browning machine gun, 7.62mm FN MAG, or 40mm Mk 19 grenade launcher. In former Eastern bloc nations, the KPV, PKT and NSV machine guns are common options.
In "open top" mounts the gunner sticks out of the vehicle and operates a gun on a pintle or ring mount. Ring mounts allow the gun to traverse 360 degrees, a pintle mount has a limited field of fire. It can be preferable to an enclosed gunner because it allows a greater field of view and communication using shouts and hand signals. However, the gunner is poorly protected and at risk of injury in the event of vehicle rollover. During the Vietnam War, M113 gunners often suffered heavy casualties.
Enclosed vehicles are equipped with turrets that allow the crew to operate the weapons system while protected by the vehicle's armour. The Soviet BTR-60 has an enclosed turret mounted with a KPV heavy machine gun with a PKT coaxial machine gun. The U.S. Assault Amphibious Vehicle, Personnel (AAVP7's) machine guns (an M2, .50 caliber MG and a Mk 19 grenade launcher) are in fully enclosed turrets (turrets typically have optics which make them more accurate).
More recently, APCs have been equipped with remote weapon systems. The baseline Stryker carries an M2 on a Protector remote weapons system.
Medical use
APCs may be used as armoured ambulances, to conduct evacuations of wounded personnel. These vehicles are equipped with stretchers and medical supplies.
According to article 19 of the Geneva Conventions, "mobile medical units of the Medical Service may in no circumstances be attacked, but shall at all times be respected and protected by the Parties to the conflict". Although article 22 allows them to carry defensive weaponry, they are typically unarmed. Under Article 39, the emblem of the medical service "shall be displayed ... on all equipment employed in the Medical Service." As such, armoured ambulances are marked with International Committee of the Red Cross (ICRC) recognized symbols.
Variants
Infantry fighting vehicle
The infantry fighting vehicle is a derivative of the APC. Various classes of infantry fighting vehicles may be deployed alongside tanks and APCs, in armoured and mechanized forces. The fundamental difference between an APC and IFV is the role they are designed for. The CFE treaty stipulates an infantry fighting vehicle is an APC with a cannon in excess of 20 mm, and with this additional firepower the vehicle is more involved in combat, providing fire support to dismounted infantry.
Mine-resistant ambush protected vehicle
Mine-Resistant Ambush Protected Vehicle (MRAPV), also known as MRAP Vehicle, is a type of armoured personnel carrier that are designed specifically to withstand land mines, improvised explosive device (IED) attacks and ambushes to save troops' lives.
Infantry mobility vehicle
"Infantry mobility vehicle" (IMV) is a new name for the old concept of an armoured car, with an emphasis on mine resistance. They are primarily used to protect passengers in unconventional warfare.
The South African Casspir was first built in the late 1970s. In the 21st century, they gained favour in the post-Cold-War geopolitical climate. Identical to earlier High Mobility Multipurpose Wheeled Vehicle (HMMWV) in design and function, the uparmoured M1114 HMMWV is a clear example of this. The addition of armour provides protection to passengers. M1114s have been largely replaced by purpose-built Mine Resistant Ambush Protected (MRAP) vehicles.
IMVs generally feature a v-shaped underbelly designed to deflect mine blasts outwards, with additional crew protection features such as four-point seat belts, and seats suspended from the roof or sides of the vehicle. Many feature a remote weapon system. Usually four-wheel drive, these IMVs are distinct from 8-, 6-, and 4-wheeled APCs (such as the VAB), being closer in appearance to civilian armoured money and gold transporters.
| Technology | Maneuver | null |
19278728 | https://en.wikipedia.org/wiki/Architectural%20engineering | Architectural engineering | Architectural engineering or architecture engineering, also known as building engineering, is a discipline that deals with the engineering and construction of buildings, such as environmental, structural, mechanical, electrical, computational, embeddable, and other research domains. It is related to Architecture, Mechatronics Engineering, Computer Engineering, Aerospace Engineering, and Civil Engineering, but distinguished from Interior Design and Architectural Design as an art and science of designing infrastructure through these various engineering disciplines, from which properly align with many related surrounding engineering advancements.
From reduction of greenhouse gas emissions to the construction of resilient buildings, architectural engineers are at the forefront of addressing several major challenges of the 21st century. They apply the latest scientific knowledge and technologies to the design of buildings. Architectural engineering as a relatively new licensed profession emerged in the 20th century as a result of the rapid technological developments. Architectural engineers are at the forefront of two major historical opportunities that today's world is immersed in: (1) that of rapidly advancing computer-technology, and (2) the parallel revolution of environmental sustainability.
Architects and architectural engineers both play crucial roles in building design and construction, but they focus on different aspects. Architectural engineers specialize in the technical and structural aspects, ensuring buildings are safe, efficient, and sustainable. Their education blends architecture with engineering, focusing on structural integrity, mechanical systems, and energy efficiency. They design and analyze building systems, conduct feasibility studies, and collaborate with architects to integrate technical requirements into the overall design. Architects, on the other hand, emphasize the aesthetic, functional, and spatial elements, developing design concepts and detailed plans to meet client needs and comply with regulations. Their education focuses on design theory, history, and artistic aspects, and they oversee the construction process to ensure the design is correctly implemented.
Subdisciplines of architectural engineering
Mechanical, electrical, and plumbing (MEP)
Mechanical engineering and electrical engineering engineers are specialists when engaged in the building design fields. This is known as mechanical, electrical, and plumbing (MEP) throughout the United States, or building services engineering in the United Kingdom, Canada, and Australia. Mechanical engineers often design and oversee the heating, ventilation and air conditioning (HVAC), plumbing, and rainwater systems. Plumbing designers often include design specifications for simple active fire protection systems, but for more complicated projects, fire protection engineers are often separately retained. Electrical engineers are responsible for the building's power distribution, telecommunication, fire alarm, signalization, lightning protection and control systems, as well as lighting systems.
Structural Engineering
Structural engineering involves the analysis and design of the built environment (buildings, bridges, equipment supports, towers and walls). Those concentrating on buildings are sometimes informally referred to as "building engineers". Structural engineers require expertise in strength of materials, structural analysis, and in predicting structural load such as from weight of the building, occupants and contents, and extreme events such as wind, rain, ice, and seismic design of structures which is referred to as earthquake engineering. Architectural engineers sometimes incorporate structural as one aspect of their designs; the structural discipline when practiced as a specialty works closely with architects and other engineering specialists.
Sustainable Engineering
Sustainable engineering involves designing or operating systems to use energy and resources in a way that maintains environmental balance and ensures that future generations can meet their own needs without compromising the natural environment. Architectural engineers are influenced by sustainable engineering principles in their education, training, and practice, integrating sustainable design strategies to create buildings and structures that minimize environmental impact and enhance energy efficiency.
Building Envelope Engineering
Building enclosure and façade engineering involves the design and management of the outer shell of a building, which acts as a barrier between the interior and exterior environments. This includes walls, roofs, windows, doors, and other components that collectively ensure the building is protected from external elements such as air, water, heat, light, and noise.
The building envelope plays a crucial role in maintaining indoor comfort by controlling temperature, humidity, and airflow. It also contributes to the building's energy efficiency by minimizing heat loss in the winter and heat gain in the summer. Engineers in this field work on making sure the envelope is structurally sound, aesthetically pleasing, and performs effectively to meet various functional requirements.
Fire Protection Engineering
Fire protection engineering is a subfield of building engineering focused on the design and application of systems and practices that prevent, control, and mitigate the impact of fires. This discipline aims to protect people, property, and the environment from the destructive effects of fire through a combination of preventive measures, detection systems, and response strategies.
Fire protection engineers use their expertise to analyze potential fire scenarios, model the spread of fire and smoke, and design systems that effectively protect lives and property. They collaborate with architects, builders, and safety officials to integrate fire protection measures into the overall design and operation of buildings and facilities.
Acoustical Engineering
Acoustical or acoustics engineering in building design focuses on controlling sound within and around buildings to create a comfortable and functional auditory environment. This discipline involves the study and application of principles to manage noise levels, improve sound quality, and ensure effective sound insulation.
Acoustical engineers work closely with architects, builders, and other engineers to integrate sound control measures into the overall design of a building. They use advanced modeling and simulation tools to predict how sound will behave in different spaces and employ various materials and techniques to achieve the desired acoustic performance. Their goal is to create environments that are acoustically comfortable, meeting the specific needs of the building's occupants and its intended use.
The architectural engineer (PE) in the United States
In many jurisdictions of the United States, the architectural engineer is a licensed engineering professional. Usually a graduate of an EAC/ABET-accredited architectural engineering university program preparing students to perform whole-building design in competition with architect-engineer teams; or for practice in one of structural, mechanical or electrical fields of building design, but with an appreciation of integrated architectural requirements. Although some states require a BS degree from an EAC/ABET-accredited engineering program, with no exceptions, about two thirds of the states accept BS degrees from ETAC/ABET-accredited architectural engineering technology programs to become licensed engineering professionals. Architectural engineering technology graduates, with applied engineering skills, often gain further learning with an MS degree in engineering and/or NAAB-accredited Masters of Architecture to become licensed as both an engineer and architect. This path requires the individual to pass state licensing exams in both disciplines. States handle this situation differently on experienced gained working under a licensed engineer and/or registered architect prior to taking the examinations. This education model is more in line with the educational system in the United Kingdom where an accredited MEng or MS degree in engineering for further learning is required by the Engineering Council to be registered as a Chartered Engineer. The National Council of Architectural Registration Boards (NCARB) facilitate the licensure and credentialing of architects but requirements for registration often vary between states. In the state of New Jersey, a registered architect is allowed to sit for the PE exam and a professional engineer is allowed to take the design portions of the Architectural Registration Exam (ARE), to become a registered architect.
Formal architectural engineering education, following the engineering model of earlier disciplines, developed in the late 19th century, and became widespread in the United States by the mid-20th century. With the establishment of a specific "architectural engineering" NCEES Professional Engineering registration examination in the 1990s, and first offering in April 2003, architectural engineering became recognized as a distinct engineering discipline in the United States. Up to date NCEES account allows engineers to apply to other states PE license "by comity".
In most license-regulated jurisdictions, architectural engineers are not entitled to practice architecture unless they are also licensed as architects. Practice of structural engineering in high-risk locations, e.g., due to strong earthquakes, or on specific types of higher importance buildings such as hospitals, may require separate licensing as well. Regulations and customary practice vary widely by state or city.
The architect as architectural engineer
In some countries, the practice of architecture includes planning, designing and overseeing the building's construction, and architecture, as a profession providing architectural services, is referred to as "architectural engineering". In Japan, a "first-class architect" plays the dual role of architect and building engineer, although the services of a licensed "structural design first-class architect"(構造設計一級建築士) are required for buildings over a certain scale.
In some languages, such as Korean and Arabic, "architect" is literally translated as "architectural engineer". In some countries, an "architectural engineer" (such as the ingegnere edile in Italy) is entitled to practice architecture and is often referred to as an architect. These individuals are often also structural engineers. In other countries, such as Germany, Austria, Iran, and most of the Arab countries, architecture graduates receive an engineering degree (Dipl.-Ing. – Diplom-Ingenieur).
In Spain, an "architect" has a technical university education and legal powers to carry out building structure and facility projects.
In Brazil, architects and engineers used to share the same accreditation process (Conselho Federal de Engenheiros, Arquitetos e Agrônomos (CONFEA) – Federal Council of Engineering, Architecture and Agronomy). Now the Brazilian architects and urbanists have their own accreditation process (CAU – Architecture and Urbanism Council). Besides traditional architecture design training, Brazilian architecture courses also offer complementary training in engineering disciplines such as structural, electrical, hydraulic and mechanical engineering. After graduation, architects focus in architectural planning, yet they can be responsible to the whole building, when it concerns to small buildings (except in electric wiring, where the architect autonomy is limited to systems up to 30kVA, and it has to be done by an Electrical Engineer), applied to buildings, urban environment, built cultural heritage, landscape planning, interiorscape planning and regional planning.
In Greece licensed architectural engineers are graduates from architecture faculties that belong to the Polytechnic University, obtaining an "Engineering Diploma". They graduate after 5 years of studies and are fully entitled architects once they become members of the Technical Chamber of Greece (TEE – Τεχνικό Επιμελητήριο Ελλάδος). The Technical Chamber of Greece has more than 100,000 members encompassing all the engineering disciplines as well as architecture. A prerequisite for being a member is to be licensed as a qualified engineer or architect and to be a graduate of an engineering and architecture schools of a Greek university, or of an equivalent school from abroad. The Technical Chamber of Greece is the authorized body to provide work licenses to engineers of all disciplines as well as architects, graduated in Greece or abroad. The license is awarded after examinations. The examinations take place three to four times a year. The Engineering Diploma equals a master's degree in ECTS units (300) according to the Bologna Accords.
Education
The architectural, structural, mechanical and electrical engineering branches each have well established educational requirements that are usually fulfilled by completion of a university program.
In Canada, a CEAB-accredited engineer degree is the minimum academic requirement for registration as a P.Eng (professional engineer) anywhere in Canada and the standard against which all other engineering academic qualifications are measured. A graduate of a non-CEAB-accredited program must demonstrate that his or her education is at least equivalent to that of a graduate of a CEAB-accredited program.
In Vietnam, the engineer's degree is called Bằng kỹ sư, the first degree after five years of study. The Ministry of Education of Vietnam has also issued separate regulations for the naming of degrees not in accordance with international regulation.
Architectural engineering as a single integrated field of study
Its multi-disciplinary engineering approach is what differentiates architectural engineering from architecture (the field of the architect): which is an integrated, separate and single, field of study when compared to other engineering disciplines.
Through training in and appreciation of architecture, the field seeks integration of building systems within its overall building design. Architectural engineering includes the design of building systems including heating, ventilation and air conditioning (HVAC), plumbing, fire protection, electrical, lighting, architectural acoustics, and structural systems. In some university programs, students are required to concentrate on one of the systems; in others, they can receive a generalist architectural or building engineering degree.
| Technology | Disciplines | null |
19282236 | https://en.wikipedia.org/wiki/Chloroflexota | Chloroflexota | The Chloroflexota are a phylum of bacteria containing isolates with a diversity of phenotypes, including members that are aerobic thermophiles, which use oxygen and grow well in high temperatures; anoxygenic phototrophs, which use light for photosynthesis (green non-sulfur bacteria); and anaerobic halorespirers, which uses halogenated organics (such as the toxic chlorinated ethenes and polychlorinated biphenyls) as electron acceptors.
The members of the phylum Chloroflexota are monoderms (that is, have one cell membrane with no outer membrane), but they stain mostly gram-negative. Many well-studied phyla of bacteria are diderms and stain gram-negative, whereas well-known monoderms that stain Gram-positive include Firmicutes (or Bacillota) (low G+C gram-positives), Actinomycetota (high-G+C gram-positives) and Deinococcota (gram-positive diderms with thick peptidoglycan).
History
The taxon name was created in the 2001 edition of Volume 1 of Bergey's Manual of Systematic Bacteriology and is the Latin plural of the name Chloroflexus, the name of the type genus of the phylum, a common practice.
In 1987, Carl Woese, regarded as one of the forerunner of the molecular phylogeny revolution, divided Eubacteria into 11 divisions based on 16S ribosomal RNA (SSU) sequences and grouped the genera Chloroflexus, Herpetosiphon and Thermomicrobium into the "green non-sulfur bacteria and relatives", which was temporarily renamed as "Chloroflexi" in Volume One of Bergey's Manual of Systematic Bacteriology.
Chloroflexota being a deep branching phylum (see Bacterial phyla), it was considered in Volume One of Bergey's Manual of Systematic Bacteriology to include a single class with the same name. Since 2001, however, new classes have been created thanks to newly discovered species, and the phylum Chloroflexi is now divided into several classes.
"Dehalococcoidetes" is a placeholder name given by Hugenholtz & Stackebrandt, 2004, after "Dehalococcoides ethenogenes" a species partially described in 1997. The first species fully described was Dehalogenimonas lykanthroporepellens, by Moe et al. 2009, but in the description of that species the class was not made official nor were families or orders laid out as the two species share only 90% 16S ribosomal RNA identity, meaning that they could fall in different families or even orders.
Recent phylogenetic analysis of the Chloroflexota has found very weak support for the grouping together of the different classes currently part of the phylum. The six classes that make up the phylum did not consistently form a well-supported clade in phylogenetic trees based on concatenated sequences for large datasets of proteins, and no conserved signature indels were identified that were uniquely shared by the entire phylum. However, the classes Chloroflexi and Thermomicrobia were found to group together consistently by both the usual phylogenetic means and the identification of shared conserved signature indels in the 50S ribosomal protein L19 and the enzyme UDP-glucose 4-epimerase. It has been suggested that the phylum Chloroflexi sensu stricto should comprise only the classes Chloroflexi and Thermomicrobia, and the other four classes ("Dehalococcoidetes," Anaerolineae, Caldilineae and Ktedonobacteria) may represent one or more independent phyla branching in the neighborhood of the Chloroflexi.
Phylogeny
The currently accepted taxonomy is based on the List of Prokaryotic names with Standing in Nomenclature (LPSN) and National Center for Biotechnology Information (NCBI).
Taxonomy
Genus "Candidatus Caldibacter" corrig. Spieck et al. 2020
Genus "Candidatus Chlorotrichoides" corrig. Oren et al. 2020 ["Candidatus Chlorothrix" Klappenbach & Pierson 2004 non Dyar 1921 non Berger-Perrot 1982]
Genus "Candidatus Nitrocaldera" Spieck et al. 2020
Genus "Candidatus Nitrotheca" Spieck et al. 2020
Class "Bathosphaeria" Mehrshad et al. 2018
Order "Bathosphaerales" Mehrshad et al. 2018
Family "Bathosphaeraceae" Mehrshad et al. 2018
Genus ?"Candidatus Bathosphaera" Mehrshad et al. 2018 (JG30-KF-CM66)
Class "Martimicrobia" Williams et al. 2024
Order "Martimicrobiales" Williams et al. 2024
Family "Martimicrobiaceae" Williams et al. 2024
Genus ?"Candidatus Laranimicrobium" Williams et al. 2024
Genus ?"Candidatus Martimicrobium" Williams et al. 2024
Class "Poriflexia" Mehrshad et al. 2018
Genus ?"Candidatus Poriflexus" Kogawa et al. 2022
Class "Spiritibacteria" Williams et al. 2024
Order "Spiritibacterales" Williams et al. 2024
Family "Spiritibacteraceae" Williams et al. 2024
Genus ?"Candidatus Aglaurobacter" Williams et al. 2024
Genus ?"Candidatus Otrerea" Williams et al. 2024
Genus ?"Candidatus Spiritibacter" Williams et al. 2024
Class "Tarhunnaeia" Williams et al. 2024
Order "Tarhunnaeales" Williams et al. 2024
Family "Tarhunnaeaceae" Williams et al. 2024
Genus ?"Candidatus Sutekhia" Williams et al. 2024
Genus ?"Candidatus Tarhunnaea" Williams et al. 2024
Class "Uliximicrobia" Williams et al. 2024
Order "Uliximicrobiales" Williams et al. 2024
Family "Uliximicrobiaceae" Williams et al. 2024
Genus ?"Candidatus Uliximicrobium" Williams et al. 2024
Class "Umbricyclopia" Mehrshad et al. 2018
Order "Umbricyclopales" Mehrshad et al. 2018
Family "Umbricyclopaceae" Mehrshad et al. 2018
Genus ?"Candidatus Umbricyclops" Mehrshad et al. 2018 (TK10)
Class "Limnocylindria" Mehrshad et al. 2018
Order "Limnocylindrales" Mehrshad et al. 2018
Family "Limnocylindraceae" Mehrshad et al. 2018 (SL56)
Genus "Candidatus Aquidulcis" corrig. Rodriguez-R et al. 2020
Genus "Candidatus Limnocylindrus" Mehrshad et al. 2018
Class Ktedonobacteria Cavaletti et al. 2007 emend. Yabe et al. 2010
Order Thermogemmatisporales Yabe et al. 2011
Family Thermogemmatisporaceae Yabe et al. 2011
Genus Thermogemmatispora Yabe et al. 2011
Order Ktedonobacterales Cavaletti et al. 2007
Family Dictyobacteraceae Wang et al. 2019
Genus Dictyobacter Yabe et al. 2017
Genus Tengunoibacter Wang et al. 2019
Family Ktedonobacteriaceae Cavaletti et al. 2007
Genus Ktedonobacter corrig. Cavaletti et al. 2007
Genus Ktedonospora Yabe et al. 2021
Family Ktedonosporobacteraceae Yan et al. 2020
Genus Ktedonosporobacter Yan et al. 2020
Family Reticulibacteraceae Yabe et al. 2021
Genus ?Reticulibacter Yabe et al. 2021
Family Thermosporotrichaceae Yabe et al. 2010
Genus Thermosporothrix Yabe et al. 2010
Class Thermomicrobiia Oren, Parte & Garrity 2016
Order Sphaerobacterales Stackebrandt, Rainey & Ward-Rainey 1997
Family Sphaerobacteraceae Stackebrandt, Rainey & Ward-Rainey 1997
Genus Nitrolancea Sorokin et al. 2014 ["Nitrolancetus" Sorokin et al. 2012]
Genus Sphaerobacter Demharter et al. 1989
Order Thermomicrobiales Garrity & Holt 2002
Family Thermomicrobiaceae Garrity & Holt 2002
Genus Thermalbibacter Zhao et al. 2023
Genus Thermomicrobium Jackson, Ramaley & Meinschein 1973
Genus Thermorudis King & King 2014
Class Chloroflexia Gupta et al. 2013
Genus ?"Candidatus Chlorohelix" Tsuji et al. 2024
Genus ?"Dehalobium" Wu et al. 2002
Genus ?"Candidatus Lithoflexus" Saghai et al. 2020
Genus ?"Candidatus Sarcinithrix" Nierychlo et al. 2019
Order "Thermobaculales" Chuvochina et al. 2023
Family "Thermobaculaceae" Chuvochina et al. 2023
Genus "Thermobaculum" Botero et al. 2004
Order Kallotenuales Cole et al. 2013
Family Kallotenuaceae Cole et al. 2013
Genus Kallotenue Cole et al. 2013
Order Herpetosiphonales Gupta et al. 2013
Family Herpetosiphonaceae Gupta et al. 2013
Genus "Candidatus Anthektikosiphon" Ward, Fischer & McGlynn 2020
Genus Herpetosiphon Holt & Lewin 1968 [Flavilitoribacter García-López et al. 2020]
Order Chloroflexales Gupta et al. 2013
Suborder Roseiflexineae Gupta et al. 2013
Family Roseiflexaceae Gupta et al. 2013 ["Kouleotrichaceae" Mehrshad et al. 2018]
Genus ?Heliothrix Pierson et al. 1986
Genus "Kouleothrix" Kohno et al. 2002
Genus "Candidatus Ribeiella" Petriglieri et al. 2023
Genus Roseiflexus Hanada et al. 2002
Suborder Chloroflexineae Gupta et al. 2013
Family Chloroflexaceae Gupta et al. 2013
Genus ?"Candidatus Chloranaerofilum" Thiel et al. 2016
Genus Chloroflexus Pierson & Castenholz 1974 ["Chlorocrinis"]
Family Oscillochloridaceae Gupta et al. 2013
Genus ?Chloronema ♪ Dubinina & Gorlenko 1975
Genus "Candidatus Chloroploca" Gorlenko et al. 2014
Genus Oscillochloris Gorlenko & Pivovarova 1989
Genus "Candidatus Viridilinea" Grouzdev et al. 2018
Class Tepidiformia Kochetkova et al. 2020
Order Tepidiformales Kochetkova et al. 2020
Family Tepidiformaceae Kochetkova et al. 2020
Genus "Candidatus Amarobacillus" Petriglieri et al. 2023
Genus ?"Candidatus Amarobacter" Petriglieri et al. 2023
Genus "Tepidiforma Kochetkova et al. 2020
Class Dehalococcoidia Löffler et al. 2013
Order Dehalococcoidales Löffler et al. 2013
Family Dehalococcoidaceae Löffler et al. 2013
Genus Dehalococcoides Löffler et al. 2013
Genus Dehalogenimonas Moe et al. 2009
Order "Lucifugimonadales" Lim et al. 2023
Family "Lucifugimonadaceae" Lim et al. 2023
Genus "Candidatus Lucifugimonas" Lim et al. 2023 [SAR202; UBA1151]
Order "Australimonadales" Prabhu et al. 2024
Family "Australimonadaceae" Prabhu et al. 2024
Genus "Candidatus Australimonas" Prabhu et al. 2024 [UBA2963]
Order "Monstramariales" Landry et al. 2017 [SAR202 group III; UBA3495]
Family "Monstramariaceae" Landry et al. 2017
Genus "Carboxydicoccus" Dede et al. 2024 [UBA11650]
Class Ardenticatenia Kawaichi et al. 2013
Order Ardenticatenales Kawaichi et al. 2013
Family Ardenticatenaceae Kawaichi et al. 2013
Genus Ardenticatena Kawaichi et al. 2013
Order "Epilineales" Petriglieri et al. 2023
Family "Epilineaceae" Petriglieri et al. 2023
Genus "Candidatus Avedoeria" Petriglieri et al. 2023
Genus "Candidatus Epilinea" Petriglieri et al. 2023
Class "Caldilineia" Oren, Parte & Garrity 2016 ex Cavalier-Smith 2020
Order Caldilineales Yamada et al. 2006
Family "Amarolineaceae" Andersen et al. 2019
Genus "Candidatus Amarolinea" Andersen et al. 2019
Family Caldilineaceae Yamada et al. 2006
Genus Caldilinea Sekiguchi et al. 2003
Genus "Candidatus Fredericiella" Petriglieri et al. 2023
Genus Litorilinea Kale et al. 2013
Class Thermoflexia Dodsworth et al. 2014
Order Thermoflexales Dodsworth et al. 2014
Family "Roseilineaceae" Ward et al. 2021
Genus "Candidatus Brachythrix" Petriglieri et al. 2023
Genus "Candidatus Roseilinea" Thiel et al. 2016
Family Thermoflexaceae Dodsworth et al. 2014
Genus Thermoflexus Dodsworth et al. 2014
Class "Thermofontia" corrig. Ward et al. 2018
Order "Phototrophicales" Zheng et al. 2022
Family "Phototrophicaceae" Zheng et al. 2022
Genus "Candidatus Flexicrinis" Petriglieri et al. 2023
Genus "Candidatus Flexifilum" Petriglieri et al. 2023
Genus "Phototrophicus" Zheng et al. 2022
Class "Anaerolineia" Oren, Parte & Garrity 2016
Genus ?"Candidatus Defluviifilum" Nierychlo et al. 2019
Family "Profundisolitariaceae" Mehrshad et al. 2018
Genus ?"Candidatus Profundisolitarius" Mehrshad et al. 2018 (CL500-11)
Order "Promineifilales" Chuvochina et al. 2023
Family "Promineifilaceae" Chuvochina et al. 2023
Genus "Candidatus Hadersleviella" Petriglieri et al. 2023
Genus "Candidatus Leptofilum" Petriglieri et al. 2023
Genus "Candidatus Leptovillus" Petriglieri et al. 2023
Genus "Candidatus Promineifilum" corrig. McIlroy et al. 2016
Genus "Candidatus Trichofilum" Petriglieri et al. 2023
Order Aggregatilineales Nakahara et al. 2019
Family Aggregatilineaceae Nakahara et al. 2019
Genus Aggregatilinea Nakahara et al. 2019
Order Anaerolineales Yamada et al. 2006
Family "Villigracilaceae" Petriglieri et al. 2023
Genus "Candidatus Defluviilinea" Petriglieri et al. 2023
Genus "Candidatus Denitrolinea" Petriglieri et al. 2023
Genus "Desulfolinea" Van Vliet et al. 2020
Genus ?"Candidatus Manresella" Petriglieri et al. 2023
Genus "Candidatus Villigracilis" Nierychlo et al. 2019 ex Petriglieri et al. 2023
Family Anaerolineaceae Yamada et al. 2006
Genus Anaerolinea Sekiguchi et al. 2003 emend. Yamada et al. 2006
Genus Bellilinea Yamada et al. 2007
Genus "Brevefilum" McIlroy et al. 2017
Genus "Candidatus Denitrolinea" Okubo et al. 2021
Genus Flexilinea Sun et al. 2015
Genus Levilinea Yamada et al. 2006
Genus Leptolinea Yamada et al. 2006
Genus Longilinea Yamada et al. 2007
Genus ?"Candidatus Mesolinea" Bedoya-Urrego & Alzate 2024
Genus Ornatilinea Podosokorskaya et al. 2013
Genus Pelolinea Imachi et al. 2014
Genus ?Thermomarinilinea Nunoura et al. 2013
Genus "Thermanaerothrix" Gregoire et al. 2011
Etymology
The name Chloroflexi is a Neolatin nominative case masculine plural of Chloroflexus, which is the name of the first genus described. The noun is a combination of the Greek adjective chloros, -a, on (χλωρός, -ά, -όν), meaning "greenish-yellow," and the Latin masculine passive perfect participle flexus (of flecto), meaning "bent." The etymology is unrelated to chlorine, an element that was discovered in 1810 by Sir Humphry Davy and named after its pale green colour. Another phylum with the same root is Chlorobiota, whereas "Cyanobacteria" has the root cyanos (κύανος), meaning "blue-green."
Unlike some other phyla, there is no theme root in the name of genera of Chloroflexota, and in fact many genera beginning with "Chloro-" or ending in "-chloris" are either cyanobacteria or chlorobi.
| Biology and health sciences | Gram-negative bacteria | Plants |
1782583 | https://en.wikipedia.org/wiki/Champsosaurus | Champsosaurus | Champsosaurus is an extinct genus of crocodile-like choristodere reptile, known from the Late Cretaceous and early Paleogene periods of North America and Europe (Campanian–Paleocene). The name Champsosaurus is thought to come from , () said in an Ancient Greek source to be an Egyptian word for "crocodiles", and , () Greek for "lizard". The morphology of Champsosaurus resembles that of gharials, with a long, elongated snout. It was native to freshwater environments where it likely preyed on fish, similar to living gharials.
History of research
Champsosaurus was the first member of the Choristodera to be described. Champsosaurus was named by Edward Drinker Cope in 1876, from isolated vertebrae found in Late Cretaceous strata of the Judith River Formation on the banks of the Judith River in Fergus County, Montana. Cope designated C. annectens as the type species rather than the first named C. profundus due to the larger number of vertebrae he attributed to the species. C. annectens was based on 9 isolated vertebral centra (AMNH FR 5696) that were not figured in the paper of which two are now lost. Cope named several other species between 1876 and 1882, also based on isolated vertebrae. Barnum Brown in 1905 described the first complete remains of Champsosaurus, and noted that one of the species attributed to Champsosaurus by Cope in 1876, C. vaccinsulensis actually represented indeterminate plesiosaur remains, and that the vertebrae that Cope used to diagnose his species of Champsosaurus were heavily eroded and the diagnostic features varied substantially along the spinal column, and were not diagnostic to species level, including the remains that Cope attributed to the type species C. annectens. The conclusion that C. annectens was undiagnostic was supported by William Parks in 1933.
Brown in 1905 named two species of Champsosaurus. One was C. ambulator, named from the specimen AMNH 983, a fragmentary skeleton with a partial skull found in the Hell Creek Formation of Montana. The other was C. laramiensis, named from AMNH 982, a nearly complete skeleton and skull, also found in the Hell Creek Formation. Parks in 1927 named C. albertensis from ROM 806, a partial skeleton lacking the skull, found in the Horseshoe Canyon Formation in Alberta, Canada. Parks in 1933 named the species C. natator from an incomplete skeleton with a fragmentary skull (TMP 81.47.1) found in the Belly River Group in the Red Deer River valley in Alberta. In 1979, Denise Sigogneau-Russell named the species C. dolloi from remains found in the Paleocene of Belgium. In 1972, Bruce Erickson named the species C. gigas from SMM P71.2.1, a partial skeleton and skull found in the Sentinel Butte Formation, Golden Valley County, North Dakota. Erickson subsequently in 1981 named the species C. tenuis from SMM P79.14.1, a partial skeleton and skull found in the Bullion Creek Formation, North Dakota. In 1998 K. Q. Gao and Richard Carr Fox described the species C. lindoei from UALVP 931, a nearly complete skeleton with skull and jaws from the Dinosaur Park Formation in Alberta. The publication also thoroughly reviewed Champsosaurus, rediagnosing most species except for C. ambulator and C. laramiensis.
Fossils of Champsosaurus have been found in North America (Alberta, Saskatchewan, Montana, New Mexico, Texas, Colorado, and Wyoming) and Europe (Belgium and France), dating from the Upper Cretaceous to the late Paleocene. Remains tentatively referred to Champsosaurus are known from the high Canadian Arctic, dating to the Coniacian–Turonian, a time of extreme warmth.
Taxonomy
Sixteen species of Champsosaurus have been named, of which seven are presently considered valid. The type species Champsosaurus annectens Cope, 1876 is considered to be dubious. The only named European species C. dolloi Sigogneau-Russell 1979 was considered to be too fragmentary to warrant a new species by Gao and Fox in 1998.
Description
Most species grew to about 1.50 m (5 ft) long, though Champsosaurus gigas, the largest species, reached 3–3.5 m (10–12 ft) in length.
Anatomy
The skull of Champsosaurus is dorsoventrally flattened, while the temporal arches are expanded posteriorly (towards the back of the skull) and laterally (away from the midline), giving the skull a heart shaped appearance when viewed from above. The snout is greatly elongated and gharial-like, making up around half the length of the skull, and at least four times as long as it is wide, with the opening of the nostrils at the end of the snout. The openings of the ears are located on the underside of the skull. The body is flat and streamlined, with heavy gastralia (rib-like bones situated in the belly). Compared to other choristoderes, the lacrimal bone is reduced in size to a small triangle, the postorbital bone does not form part of the orbit (eye socket), there is no contact between the premaxilla and the vomer bones, an internarial bone is present, the choanae are located posteriorly in correlation to the elongation of the vomer, the interpterygoid vacuity is small and completely enclosed by the pterygoid bones and located near the posterior margin of the suborbital fenestra, the shape of the suborbital fenestra is shortened and kidney like, the articulation between the pterygoid and the parasphenoid is fused, the joint between the skull and the lower jaws is anterior to level of the occipital condyles, the neomorphic bone forms most of the border of the posttemporal fenestra, the paroccipital process is strongly deflected downwards, the basal tubera of the basisphenoid are wing-like in shape and expanded backwards and downwards, the mandibular symphysis (connection of the two halves of the lower jaw) is elongated to over half the length of the tooth row, and the splenial bone strongly intervenes in the mandibular symphysis.
Internal cranial anatomy
The braincase of Champsosaurus is poorly ossified at the front of the skull (anterior), but is well ossified in the rear (posterior) similar to other diapsids. The cranial endocast (space occupied by the brain in the cranial vault) is similar to that of basal archosauromorphs, being proportionally narrow in both dorsoventral and lateral axes, with an enlarged pineal body and olfactory bulbs. The optic lobes and flocculi are small in size, indicating only average vision ability at best. The olfactory chambers of the nasal passages and olfactory stalks of the braincase are reasonably large, indicating that Champsosaurus probably had good olfactory capabilities (sense of smell). The nasal passages lack bony turbinates. The semicircular canals are most similar to those of other aquatic reptiles. The expansion of the sacculus indicates that Champsosaurus likely had an increased sensitivity to low frequency sounds and vibrations. The absence of an otic notch indicates that Champsosaurus lacked a tympanum, and probably had a poor ability to detect airborne sounds.
Teeth
Champsosaurus, like many of its fellow neochoristoderes, features teeth with striated enamel of the tooth crown with enamel infolding at the base. Anterior teeth are typically sharper and more slender than posterior teeth. Like other choristoderes, Champsosaurus possessed palatal teeth (teeth present on the bones of the roof of the mouth), with longitudinal rows present on the pterygoid, palatine and vomer, alongside a small row on the flange of the pterygoid. The palatine teeth of Champsosaurus are located on raised platforms of bone, though the wideness of the platforms, the sharpness and orientation of teeth vary between species. The orientation of the teeth varies in the jaw, with the posterior teeth being orientated backward. The palatal teeth, likely in combination with a fleshy tongue probably aided in gripping and swallowing prey.
Skin
Skin impressions of Champsosaurus have been reported. They consist of small (0.6-0.1 mm) pustulate and rhomboid scales, with the largest scales being located on the lateral sides of the body, decreasing in size dorsally, no osteoderms were present.
Classification
Champsosaurus belongs to the Neochoristodera, a clade within Choristodera, the members of which are characterised by elongated snouts and expanded temporal arches. The group first appeared during the Early Cretaceous in Asia, and are suggested to have evolved in the regional absence of aquatic crocodyliformes. While Neochoristodera is a well supported grouping, the relationships of the members of the group to each other are uncertain, with the clade having been recovered as a polytomy in recent analyses.
Phylogeny of Choristodera after Dong and colleagues (2020).
Paleobiology
Champsosaurus is thought to have been highly specialised for aquatic life. Erickson 1985 suggested that the expanded temporal arches, which likely anchored powerful jaw muscles, and elongated snout allowed Champsosaurus to prey on fish akin to modern gharials, with these adaptions allowing rapid movement of the head and jaws for prey capture. A study in 2021 found that the middle and posterior neck vertebrae of Champsosaurus were adapted for lateral movement, and that Champsosaurus may have fed by laterally sweeping its head, using its slender jaws to grab individual fish from shoals, akin to how modern gharials feed. The mechanism of head movement is different from that of gharials, where the lateral movement occurs at the head-neck joint. It is unlikely that Champsosaurus fed by inertial feeding (where the prey is temporarily let go and the head moved forwards in order to force the prey deeper into the throat), but that the prey was moved down the throat by the tongue in combination with the palatal dentition. Erickson 1985 proposed that the position of the nostrils at the front of the snout allowed Champsosaurus to spend large amounts of time at the bottom of water bodies, with the head being angled upwards to allow the snout to act like a snorkel when the animal needed to breathe. However, later studies suggested that the neck vertebrae of Champsosaurus only had a limited ability to flex upwards. Champsosaurus co-existed with similarly sized aquatic crocodilians and at some Paleocene localities with fellow neochoristodere Simoedosaurus, though in assemblages where Champsosaurus occurs longirostrine (long snouted) gharial-like crocodilians are absent, suggesting that there was niche differentiation. Previously, two species of Champsosaurus were identified from the Tullock Formation in Montana. However, these differences are now thought to be sexually dimorphic, with presumed females possessing robust limb bones. Non-deformation related fusion of the sacral vertebrae is also observed in specimens with robust limb bones. These are hypothesised to be related to breeding behaviour, with the more robust limb bones and fused sacrals of the females allowing them to move themselves onto land to lay eggs.
| Biology and health sciences | Choristoderes | Animals |
1782897 | https://en.wikipedia.org/wiki/SN%201054 | SN 1054 | SN 1054 is a supernova that was first observed on , and remained visible until .
The event was recorded in contemporary Chinese astronomy, and references to it are also found in a later (13th-century) Japanese document and in a document from the Islamic world. Furthermore, there are a number of proposed references from European sources recorded in the 15th century, as well as a pictograph associated with the Ancestral Puebloan culture found near the Peñasco Blanco site in New Mexico, United States. The pyramids at Cahokia in the midwestern United States may have been built in response to the supernova's appearance in the sky.
The remnant of SN 1054, which consists of debris ejected during the explosion, is known as the Crab Nebula. It is located in the sky near the star Zeta Tauri (ζ Tauri). The core of the exploding star formed a pulsar, called the Crab Pulsar (or PSR B0531+21). The nebula and the pulsar that it contains are some of the most studied astronomical objects outside the Solar System. It is one of the few Galactic supernovae where the date of the explosion is well known. The two objects are the most luminous in their respective categories. For these reasons, and because of the important role it has repeatedly played in the modern era, SN 1054 is one of the best known supernovae in the history of astronomy.
The Crab Nebula is easily observed by amateur astronomers thanks to its brightness, and was also catalogued early on by professional astronomers, long before its true nature was understood and identified. When the French astronomer Charles Messier watched for the return of Halley's Comet in 1758, he confused the nebula for the comet, as he was unaware of the former's existence. Motivated by this error, he created his catalogue of non-cometary nebulous objects, the Messier Catalogue, to avoid such mistakes in the future. The nebula is catalogued as the first Messier object, or M1.
Identification of the supernova
The Crab Nebula was identified as the supernova remnant of SN 1054 between 1921 and 1942, at first speculatively (1920s), with some plausibility by 1939, and beyond reasonable doubt by Jan Oort in 1942.
In 1921, Carl Otto Lampland was the first to announce that he had seen changes in the structure of the Crab Nebula. This announcement occurred at a time when the nature of the nebulae in the sky was completely unknown. Their nature, size and distance were subject to debate. Observing changes in such objects allows astronomers to determine whether their spatial extension is "small" or "large", in the sense that notable fluctuations to an object as vast as our Milky Way cannot be seen over a small time period, such as a few years, whereas such substantial changes are possible if the size of the object does not exceed a diameter of a few light-years. Lampland's comments were confirmed some weeks later by John Charles Duncan, an astronomer at the Mount Wilson Observatory. He benefited from photographic material obtained with equipment and emulsions that had not changed since 1909; as a result the comparison with older snapshots was easy and emphasized a general expansion of the cloud. The points were moving away from the centre, and did so faster as they got further from it.
Also in 1921, Knut Lundmark compiled the data for the "guest stars" mentioned in the Chinese chronicles known in the West. He based this on older works, having analysed various sources such as the Wenxian Tongkao, studied for the first time from an astronomical perspective by Jean-Baptiste Biot in the middle of the 19th century. Lundmark gives a list of 60 suspected novae, then the generic term for a stellar explosion, in fact covering what is now understood as two distinct phenomena, novae and supernovae. The nova of 1054, already mentioned by the Biots in 1843, is part of the list. It stipulates the location of this guest star in a note at the bottom of the page as being "close to NGC 1952", one of the names for the Crab Nebula, but it does not seem to create an explicit link between them.
In 1928, Edwin Hubble was the first to note that the changing aspect of the Crab Nebula, which was growing bigger in size, suggests that it is the remains of a stellar explosion. He realised that the apparent speed of change in its size signifies that the explosion which it comes from occurred only nine centuries ago (as observed on Earth), which puts the date of the explosion in the period covered by Lundmark's compilation. He also noted that the only possible nova in the region of Taurus (where the cloud is located) is that of 1054, whose age is estimated to correspond to an explosion dating from the start of the second millennium.
Hubble therefore deduced, correctly, that this cloud was the remains of the explosion which was observed by Chinese astronomers.
Hubble's comment remained relatively unknown as the physical phenomenon of the explosion was not known at the time. Eleven years later, when the fact that supernovae are very bright phenomena was highlighted by Walter Baade and Fritz Zwicky and when their nature was suggested by Zwicky, Nicholas Mayall proposed that the star of 1054 was actually a supernova, based on the speed of expansion of the cloud, measured by spectroscopy, which allows astronomers to determine its physical size and distance, which he estimated at 5000 light-years. This was under the assumption that the velocities of expansion along the line of sight and perpendicularly to it were identical. Based on the reference to the brightness of the star which featured in the first documents discovered in 1934, he deduced that it was a supernova rather than a nova.
This deduction was subsequently refined, which pushed Mayall and Jan Oort in 1942 to analyse historic accounts relating to the guest star more closely (see below). These new accounts, globally and mutually concordant, confirm the initial conclusions by Mayall and Oort in 1939 and the identification of the guest star of 1054 is established beyond all reasonable doubt.
Most other historical supernovas are not confirmed so conclusively: supernovas of the first millennium (SN 185, SN 386 and SN 393) are established on the basis of a single document each, and so they cannot be confirmed; in relation to the supposed historical supernova which followed the one in 1054, SN 1181, there are legitimate doubts concerning the proposed remnant (3C58) and an object of less than 1000 years of age. Other historical supernovae of which there are written accounts which precede the invention of the telescope (SN 1006, SN 1572 and SN 1604) are however established with certitude. Telescope-era supernovae are of course associated with full certitude with their remnant, when one is observed, but none is known within the Milky Way.
Historical records
SN 1054 is one of eight supernovae in the Milky Way that can be identified because written testimony describing the explosion has survived. In the nineteenth century, astronomers began to take an interest in the historic records. They compiled and examined the records as part of their research on recent novae, comets, and later, the supernovae.
The first Westerners to attempt a systematic compilation of records from China were the father and son Biot. In 1843, the sinologist Édouard Biot translated for his father, the astronomer and physicist Jean-Baptiste Biot, passages from the astronomical treatise of the 348-volume Chinese encyclopaedia, the Wenxian Tongkao.
Almost 80 years later in 1921, Knut Lundmark undertook a similar effort based on a greater number of sources. In 1942, Jan Oort, convinced that the Crab Nebula was the "guest star" of 1054 described by the Chinese, asked sinologist J.J.L. Duyvendak to help him compile new evidence on the observation of the event.
Chinese astronomy
Star-like objects that appeared temporarily in the sky were generically called "guest stars" (kè xīng 客星) by Chinese astronomers. The guest star of 1054 occurred during the reign of the Emperor Renzong of the Song dynasty (960–1279). The relevant year is recorded in Chinese documents as "the first year of the Zhihe era". Zhihe was an era name used during the reign of Emperor Renzong, and corresponds to the years 1054–1056, so the first year of the Zhihe era corresponds to the year 1054.
Some of the Chinese accounts are well preserved and detailed. The oldest and most detailed accounts are from Song Huiyao and Song Shi, historiographical works of which the extant text was redacted perhaps within a few decades of the event. There are also some later records, redacted in the 13th century, which are not necessarily independent of the older ones.
Three accounts are apparently related because they describe the angular distance from the guest star to Zeta Tauri as "perhaps several inches away", but they are in apparent disagreement about the date of appearance of the star. The older two mention the day jichou 己丑, but the third, the Xu Zizhi Tongjian Changbian, the day yichou 乙丑. These terms refer to the Chinese sexagenary cycle, corresponding to numbers 26 and 2 of the cycle, which corresponds, in the context where they are cited, respectively, to 4 July and 10 June.
As the redaction of the third source is of considerably later date (1280) and the two characters are similar, this is easily explained as a transcription error, the historical date being jichou 己丑, 4 July.
The description of the guest star's location as "to the south-east of Tianguan, perhaps several inches away" has perplexed modern astronomers, because the Crab Nebula is not situated in the south-east, but to the north-west of Zeta Tauri.
The duration of visibility is explicitly mentioned in chapter 12 of Song Shi, and slightly less accurately, in the Song Huiyao. The last sighting was on 6 April 1056, after a total period of visibility of 642 days.
This duration is supported by the Song Shi. The Song Huiyao by contrast mentions a visibility of the guest star of only 23 days, but this is after mentioning visibility during daylight. This period of 23 days applies in all likelihood solely to visibility during the day, which naturally was much shorter.
Sources
The Song Huiyao (literally "Collected important documents of the Song dynasty") covers the period 960–1220. Huiyao is a traditional form of history books in China which aimed mainly to preserve primary sources, and as such are important sources supplementing the official Twenty-Four Histories. The Song dynasty had a specific government department dedicated to compiling the Huiyao, and some 2,200 volumes were published in ten batches during the Song dynasty. However, most of these documents were lost by the time of the Qing Dynasty except for the synopsis and a relatively small portion preserved as part of the imperial Yongle Encyclopedia. In 1809, the portion preserved in the Yongle Encyclopedia was extracted and re-published as the Song Huiyao Jigao (the "draft extract of the Song Huiyao"). Subsequent scholars have worked on the project further and the current edition dates from 1936.
This document recounts the observation of the guest star, focusing on the astrological aspect but also giving important information on the visibility of the star, by day and by night.
The Song Shi is the official annals of the Song dynasty. Chapter 12 mentions the guest star, not its appearance but rather the moment of its disappearance. The corresponding entry dated 6 April 1056 indicates:
In chapter 56 ("Astronomical treaty") of the same document, the guest star is again mentioned in a chapter dedicated to this type of phenomenon, this time focusing on its appearance,
The Xu Zizhi Tongjian Changbian ("Long compilation of the continuation of the Zizhi Tongjian"), a book covering the period of 960–1126 and written 40 years or so later by Li Tao (1114–1183), contains the oldest Chinese testimonies relating to the observation of the star. It was rediscovered in 1970 by the specialist in Chinese civilisations Ho Peng Yoke and collaborators.
It is relatively imprecise in the case of the explosion of SN 1054. A loose translation of what was stated:
There is an account of the star from the Liao dynasty, which ruled in the area around northeast China from 907 to 1125. The book in question, the Qidan Guo Zhi, was compiled by Ye Longli in 1247. It includes various astronomical notes, some of which are clearly copied from the Song Shi. This entry referring to the star of 1054 seems unique:
The account of Qidan Guo Zhi alluded to the notable astronomical events that preceded the death of King Xingzong. Various historical documents allow us to establish the date of death of the Emperor Xingzong as 28 August 1055, during the eighth lunar month of the twenty-fourth (and not twenty-third) year of his reign. The dates of the two astronomical events mentioned (the eclipse and the appearance of the guest star) are not specified, but were probably before the obituary (2 or 3 years at most). Two solar eclipses were visible shortly before that date in the Khitan kingdom, on 13 November 1053 and 10 May 1054. Of these, only one occurred around noon, that of 13 November; it seems likely that this is what the document mentions. As for the guest star, only a rough estimate of location is given, corresponding to the moon mansion Mao. This mansion is situated just east of where the star appeared, as mentioned in the other testimonies. Since no other known significant astronomical event occurred in this region of the sky during the two years that preceded the death of Xingzong, it seems likely that the text is actually referring to the star of 1054.
The Wenxian Tongkao is the first East Asian source that came to the attention of Western astronomers; it was translated by Édouard Biot in 1843. This source, compiled by Ma Duanlin in 1280, is relatively brief. The text is very close to that of the Song Shi:
Identity of Tianguan
The asterisms (or "constellations") of Chinese astronomy were catalogued around the 2nd century BC. The asterisms with the brightest stars in the sky were compiled in a work called Shi Shi, which also includes Tianguan.
Identification of Tianguan is comparatively easy, as it is indicated that it is located at the foot of the Five Chariots asterism, the nature of which is left in hardly any doubt by representation on maps of the Chinese sky: it consists of a large pentagon containing the bright stars of the Auriga. As Tianguan is also represented to the north of the Three Stars asterism, the composition of which is well known, corresponding to the bright stars of Orion, its possible localisation is strongly restricted to the immediate proximity of the star ζ Tauri, located between "Five Chariots" and "Three Stars". This star, of medium brightness (apparent magnitude of 3.3), is the only star of its level of brightness in this area of the sky (there is no other star that is brighter than an apparent magnitude of 4.5 within 7 degrees of ζ Tauri), and therefore the only one likely to figure among the asterisms of "Shi Shi". All of these elements, along with some others, allow "Tianguan" to be confirmed beyond reasonable doubt as corresponding to the star ζ Tauri.
Position relative to Tianguan
Three Chinese documents indicate that the guest star was located "perhaps a few inches" South-East of Tianguan. Song Shi and Song Huiyao stipulate that it "was standing guard" for the asterism, corresponding to the star ζ Tauri. The "South-East" orientation has a simple astronomical meaning, the celestial sphere having, like the Earth's globe, both north and south celestial poles, the "South-East" direction thus corresponding to a "bottom-left" location in relation to the reference object (in this case, the star ζ Tauri) when it appears at the South. However, this "South-East" direction has long left modern astronomers perplexed in the context of this event: the logical remnant of the supernova corresponding to the guest star is the Crab Nebula, but it is not situated to the southeast of ζ Tauri, rather in the opposite direction, to the northwest.
The term "perhaps a few inches" (ke shu cun in the Latin transliteration) is relatively uncommon in Chinese astronomical documents. The first term, ke, is translated as "approximately" or "perhaps", the latter being currently preferred. The second term, shu, means "several", and more specifically any number between 3 and 9 (limits included). Finally, cun resembles a unit of measurement for angles translated by the term "inch".
It is part of a group of three angular units, zhang (also written chang), chi ("foot") and cun ("inch"). Different astronomical documents indicate without much possible discussion that a zhang corresponds to ten chi, and that one chi corresponds to ten cun. The angular units are not the ones used to determine stars' coordinates, which are given in terms of du, an angular unit corresponding to the average angular distance travelled by the sun per day, which corresponds to around 360/365.25 degrees, in other words almost one degree. The use of different angular units can be surprising, but it is similar to the current situation in modern astronomy, where the angular unit used to measure angular distances between two points is certainly the same as for declination (the degree), but is different for right ascension (which is expressed in angular hours; an angular hour corresponds to exactly 15 degrees). In Chinese astronomy, right ascension and declination have the same unit, which is not the one used for other angular distances. The reason for this choice to use different units in the Chinese world is not well known.
Meaning of units
However, the exact value of these new units (zhang, chi and cun) was never stipulated, but can be deduced by the context in which they are used. For example, the spectacular passing of Halley's comet in 837 indicates that the tail of the comet measured 8 zhang. Even if it is not possible to know the angular size of the comet at the time it passed, it is certain that 8 zhang correspond to 180 degrees at the most (maximum visible angle on the celestial sphere), which means that one zhang can hardly exceed 20 degrees, and therefore one cun cannot exceed 0.2 degrees. A more rigorous estimation was made from 1972 on the basis of references of minimal separations expressed in chi or cun between two stars in the case of various conjunctions.
The results suggest that one cun is between 0.1 and 0.2 degrees and that one chi is between 0.44 and 2.8 degrees, a range which is compatible with the estimations for one cun. A more solid estimation error is that it is generally accepted that one chi is in the order of one degree (or one du), and that one cun is in the order of one tenth of a degree. The expression "perhaps a few inches" therefore suggests an angular distance in the order of one degree or less.
Problems with description
If all the available elements strongly suggest that the star of 1054 was a supernova, and that in the area next to where the star was seen, there is a remnant of a supernova which has all of the characteristics expected of an object that is around 1,000 years old, a major problem arises: the new star is described as being to the South-East of Tianguan, while the Crab Nebula is to the North-East. This problem has been known since the 1940s and has long been unsolved. In 1972 for example, Ho Peng Yoke and his colleagues suggested that the Crab Nebula was not the product of the explosion of 1054, but that the true remnant was to the South-East, as indicated in several Chinese sources. For this, they envisaged that the angular unit cun corresponds to a considerable angle of 1 or 2 degrees, meaning that the distance from the remnant to ζ Tauri was therefore considerable. Aside from the fact that this theory does not account for the large angular sizes of certain comets, expressed in zhang, it comes up against the fact that there it does not make sense to measure the gap between a guest star and a star located so far away from it, when there are closer asterisms that could be used.
In their controversial article (see European sightings, below) Collins and his colleagues make another suggestion: on the morning of 4 July, the star ζ Tauri was not bright enough and too low on the horizon to be visible. If the guest star, which was located close to it, was visible, it is only because its brightness was comparable to Venus. However, there was another star, brighter and higher on the horizon, which was possibly visible, for reference: Beta Tauri (β Tauri). This star is located at around 8 degrees north-north-west of ζ Tauri. The Crab Nebula is south-south-east of β Tauri. Collins et al. suggest therefore that at the time of its discovery, the star was seen to the south-east of β Tauri, and that as the days passed and visibility improved, astronomers were able to see that it was in fact a lot closer to ζ Tauri, but that the direction "south-east" used for the first star was kept in error.
The solution to this problem was suggested (without proof) by A. Breen and D. McCarthy in 1995 and proved very convincingly by D. A. Green et F. R. Stephenson (2003). The term "stand on guard" obviously signifies a proximity between the two stars, but also means a general orientation: a guest star "standing on guard" for a fixed star is systematically located below it. In order to support this theory, Green and Stephenson investigated other entries in Song Shi, which also includes reference to "standing on guard". They selected entries relating to conjunctions betweens the stars identified and planets, of which the trajectory can be calculated without difficulty and with great precision on the indicated dates. Of the 18 conjunctions analysed, spreading from 1172 (the Jupiter–Regulus conjunction on 5 December) to 1245 (the Saturn–Gamma Virginis conjunction on 17 May), the planet was more to the north (in the sense of a lower declination) in 15 cases, and in the three remaining cases, it was never in the south quadrant of the star.
In addition, Stephenson and Clark (1977) had already highlighted such an inversion of direction in a planetary conjunction: on 13 September 1253, an entry in the astronomical report Koryo-sa indicated that Mars had hidden the star to the south-east of the twenty-eight mansions sign Ghost (Delta Cancri), while in reality, it approached the star north-west of the asterism (Eta Cancri).
Meigetsuki (Japan)
The oldest and most detailed record from Japan is in the Meigetsuki, the diary of Fujiwara no Teika, a poet and courtier.
There are two other Japanese documents, presumably dependent on the Meigetsuki:
The 14th century Ichidai Yoki: The description is very similar to the Meigetsuki, omitting several details (hour of apparition, and possibly erroneous parts of the lunar month). The short text also contains many typographical errors.
The 17th-century Dainihonshi, containing very little information. The brevity contrasts with the more detailed descriptions of "guest stars" (supernovas) of 1006 and 1181.
The Meigetsuki places the event in the fourth lunar moon, one month earlier than the Chinese texts.
Whatever the exact date during this month, there seems to be a contradiction between this period and the observation of the guest star: the star was close to the sun, making daytime and nighttime observation impossible. The visibility in daylight as described by the Chinese texts is thus validated by the Japanese documents, and is consistent with a period of moderate visibility, which implies that the star's period of diurnal visibility was very short.
In contrast, the day of the cycle given in the Chinese documents is compatible with the months that they state, reinforcing the idea that the month on the Japanese document is incorrect.
The study of other medieval supernovas (SN 1006 and SN 1181) reveals a proximity in the dates of discovery of a guest star in China and Japan, although clearly based on different sources.
Fujiwara no Teika's interest in the guest star seems to have come accidentally whilst observing a comet in December 1230, which prompted him to search for evidence of past guest stars, among those SN 1054 (as well as SN 1006 and SN 1181, the two other historic supernovas from the early second millennium). The entry relating to SN 1054 can be translated as:
The source used by Fujiwara no Teika is the records of Yasutoshi Abe (Onmyōdō doctor), but it seems to have been based, for all the astronomical events he has recorded, on documents of Japanese origin.
The date he gives is prior to the third part of ten days of the lunar month mentioned, which corresponds to the period of between 30 May and 8 June 1054 of the Julian calendar, which is around one month earlier than Chinese documentation. This difference is usually attributed to an error in the lunar months (fourth place and fifth place).
The location of the guest star, clearly straddling the moon mansions Shen and Zuixi, corresponds to what would be expected of a star appearing in the immediate vicinity of Tianguan.
Ibn Butlan (Iraq)
While SN 1006, which was significantly brighter than SN 1054, was mentioned by several Arab chroniclers, there exist no Arabic reports relating to the rather faint SN 1181.
Only one Arabic account has been found concerning SN 1054, whose brightness is between those of the last two stars mentioned. This account, discovered in 1978, is that of a Nestorian Christian doctor, Ibn Butlan, transcribed in the Uyun al-Anba, a book on detailed biographies of physicians in the Islamic world compiled by Ibn Abi Usaybi'a (1194–1270) in the mid-thirteenth century. This is a translation of the passage in question:
The three years cited (AH 445, 446, 447) refer, respectively, to: 23 April 1053 – 11 April 1054, 12 April 1054 – 1 April 1055, and 2 April 1055 – 20 March 1056. There is an apparent inconsistency in the year of occurrence of the star, first announced as 446, then 445. This problem is solved by reading other entries in the book, which quite explicitly specify that the Nile was low at 446.
This year of the Muslim calendar ran from 12 April 1054 to 1 April 1055, which is compatible with the appearance of the star in July 1054, as its location (admittedly rather vague), is in the astrological sign of Gemini (which, due to axial precession, covers the eastern part of the Constellation Taurus). The date of the event in 446 is harder to determine, but the reference to the level of the Nile refers to the period preceding its annual flood, which happens during the summer.
Suggested European sightings
Since 1980, several European documents have been identified as possible observations of the supernova.
The first such suggestion was made in 1980 by Umberto Dall'Olmo (1925–1980). The following passage which reports an astronomical sighting is taken from an account compiled by Jacobus Malvecius in the 15th century:
The date this passage refers to is not explicit, however, and by means of a reference to an earthquake in Brescia 11 April 1064, it would seem ten years too late. Dall'Olmo suggests this is due to a transcription error.
Another candidate is the Cronaca Rampona, proposed in 1981, which however also indicates a date several years after the event, in 1058 instead of 1054.
The European candidate documents are imprecise, especially lacking in astronomical terms likely due to European scholars having lost many of the astronomical skills of antiquity. In contrast, the Chinese accounts pin-point within a degree where the supernova occurred, as well as how long it lasted and roughly how bright it became.
The lack of accounts from European chroniclers has long raised questions. In fact, it is known that the supernova of 1006 was recorded in a large number of European documents, albeit not in astronomical terms. Among the proposed explanations for the lack of European accounts of SN 1054, its concurrence with the East-West Schism is prominent.
In fact, the date of the excommunication of the Patriarch of Constantinople Michael I Cerularius (16 July) corresponds to the star reaching its maximum brightness and being visible in the daytime.
Among the six proposed European documents, one does not seem to correspond to the year of the supernova (the chronicle of Jacobus Malvecius). Another (the Cronaca Rampona) has large dating and internal coherence problems.
The four others are relatively precisely dated, but they date from Spring and not Summer 1054, that is to say before the conjunction between the supernova and the Sun (although a Khitan document suggests this may have been possible). Three of the documents (the chronicle of Jacobus Malvecius, the Cronaca Rampona and the Armenian chronicle) make reference relatively explicitly to conjunctions between the Moon and stars, of which one is identified (Jupiter, in the Armenian chronicle).
The three other documents are very unclear.
In 1999, George W. Collins and his colleagues defended the plausibility of European sighting of SN 1054. They argue that the records suggest that European sightings even predate Chinese and Japanese reports by more than two months (April 1054). These authors emphasize the problems associated with the Chinese reports, especially the position of the supernova relative to Zeta Tauri. They also adduce a Khitan document which they suggest might establish observation of the supernova at the time of the solar eclipse of 10 May 1054 (which would corrobate the "late" date of Chinese observation of the event).
Conversely, they interpret the European documents, taken in conjunction, as plausibly establishing that an unusual astronomical phenomenon was visible in Europe in the spring of 1054, i.e. even before the Sun's conjunction with Zeta Tauri.
They also surmise that the correct year in the report by Ibn Butlan is AH 445 (23 April 1053 – 11 April 1054) rather than AH 446 (12 April 1054 – 1 April 1055).
The publication by Collins et al. was criticized by Stephenson and Green (2003). These authors insist that the problems with the Chinese and Japanese documents can easily be resolved philologically (as common copyists' mistakes) and need not indicate unreliability of the Chinese observations. Stephenson and Green condemn attempts at uncovering European sightings of the supernova as it were at any cost as suffering from confirmation bias, "anxious to ensure that this event was recorded by Europeans".
They also reject the idea of the Khitan document referring to the supernova as a mistake based in a translation of the document.
The Cronaca Rampona
The European account of a supernova sighting that is considered the most plausible is part of a medieval chronicle from the region of Bologna, the Cronaca Rampona. This text came to astronomers' attention in 1972, and was interpreted as a possible sighting of the supernova in 1981, and again in 1999. The relevant part of the chronicle translates to:
Before even looking for potential problems in the astronomical last sentence of the passage, skeptics point out at least two discrepancies in the dating: Pope Stephen IX became Pope in 1057, not 1058, and Emperor Henry III, Holy Roman Emperor was born in 1017, 39 and not 49 years before 1058, his reign having started in 1039 (as King of the Romans, then as emperor of the Romans from 1046 after Pope Clement II consecrated him during his brief pontificate). Henry III died in 1056, and his reign did not overlap with Stephen IX's papacy. It seems likely that the text underwent various alterations, as its date format uses a mix of Roman and Arabic numerals (the number 1058 is for instance written as Ml8) which was common in the 15th century when the Cronaca Rampona was assembled, but not in the 11th century when the events occurred.
Associating the stella clarissima with the 1054 supernova also requires assuming that its entry in the Cronaca Rampona is out of order, as the entries are otherwise in chronological order and the two previous entries are later than 1054 (in order, the previous entries refer to 1046, 1049, 1051, 1055, 1056, all written in a mix of Arab and Roman characters, namely Mxl6, Mxl9, Mli, Mlv and Ml6). Additionally, the date of the new moon is discrepant. Calculating the phase of the moon for every day of 1054 and converting the calends, which refer to the Roman calendar, to our Gregorian calendar shows that no month of that year had a new moon on the thirteenth day of its Calends . All of this strongly contrasts with the general precision of references to eclipse dates in medieval European chronicles: a study of 48 partial or total solar eclipses from 733 to 1544 finds that 42 dates out of 48 are correct, and of the six remaining, three are incorrect by one of two days and the three others give the correct day and month, but a wrong year.
Even assuming that the stated event nevertheless corresponds to May or June 1054, and that it describes a conjunction between the already visible supernova and the moon, a final problem arises: the moon did not pass very close to the location of the supernova during two those months.
It is therefore possible that the account instead describes an approach or a concealment of a planet by the Moon, contemporary to the date written in the document (1058). This scenario is corroborated by two perfectly dated contemporary documents which describe a conjunction and a planetary concealment by the Moon in relatively similar terms. These two documents, unearthed by Robert Russell Newton, are taken from the Annales Cavenses, Latin chronicles from la Trinità della Cava (Province of Salerno). They mention
"a bright star that entered into the circle of the (new) moon"
for both 17 February 1086 ([Martii incipiente nocte] stella clarissima in circulum lunae primae ingressa est) and for 6 August 1096 (stella clarissima venit in circulum lunae). The first event can be verified as Venus being eclipsed by the Moon, the second as the Moon passing Jupiter at a distance of less than one degree after a lunar eclipse which was also mentioned in the chronicle.
Hayton of Corycus
The Cronaca Rampona account is apparently also reflected in the Armenian chronicle of Hayton of Corycus (written before 1307).
The relevant passage translated from the Armenian manuscript reads:
Vahe Gurzadyan's proposal connecting the Hayton of Corycus's chronicle with Cronaca Rampona and SN 1054 dates to 2012.
Other
In a work entitled De Obitu Leonis ("On the Death of [Pope] Leo") by one subdeacon Libuinus, there is a report of an unusual celestial phenomenon. A certain Albertus, leading a group of pilgrims in the region of Todi, Umbria, reportedly confirmed having seen, on the day that Pope Leo IX died, a phenomenon described as
Guidoboni et al. (1994) proposed that this may relate to SN 1054, and was endorsed by Collins et al. (1999).
Guidoboni et al. (1994) also proposed a Flemish text as an account of a sighting of the supernova. The text, from Saint Paul's church—no longer extant—in the Flemish town of Oudenburg, describes the death of Pope Leo IX in Spring 1054 (the date described corresponds to 14 April 1054).
McCarthy and Breen (1997) proposed an extract from an Irish chronicle as a possible European sighting of the supernova. This chronicle indicates the following for 1054:
The date of the event corresponds to 24 April: (Saint George's Day is 23 April and fell on a Saturday in 1054. Thus the mention of the "Sunday of Saint George's Day" corresponds to the next day, 24 April) long before the sighting noted by the Chinese. The astronomical nature of the account remains very uncertain, and interpretation as a solar halo or aurora seems at least as probable.
Suggested records in North American petroglyphs
Two Native American paintings in Arizona show a crescent moon located next to a circle that could represent a star. In 1955, optical engineer and amateur archaeologist William C. Miller proposed that this represents a conjunction between the moon and the supernova, made possible by the fact that, seen from the Earth, the supernova occurred in the path of the Ecliptic. On the morning of 5 July, the moon was located in the immediate proximity of the supernova, and this proximity might have been represented in these paintings. This theory is compatible with the uncertain dating of these paintings but cannot be confirmed. The dating of the paintings is extremely imprecise (between the 10th and 12th century), and only one of them shows the crescent moon with the correct orientation in relation to the supernova on the date of the explosion. Moreover, this type of drawing could well represent a proximity of the moon with Venus or Jupiter.
Another, better known document was updated during the 1970s at the Chaco Canyon site (New Mexico), occupied around 1000 AD by the Ancestral Pueblo Peoples. On the flat underside of an overhang, it represents a hand, below which there is a crescent moon facing a star at the bottom-left. On the wall underneath the petroglyph there is a drawing which could be the core and tail of a comet. Apart from the petroglyph, which could represent the configuration of the moon and supernova on the morning of 5 July 1054, this period corresponds to the apogee of the Ancestral Pueblo civilization. It seems possible to propose an interpretation of the other petroglyph, which, if it is more recent than the other one, could possibly correspond to the passing of Halley's Comet in 1066. Although plausible, this interpretation is impossible to confirm and does not explain why it was the supernova of 1054 that was represented, rather than the supernova of 1006, which was brighter and also visible to this civilisation.
Suggested records in Aboriginal oral tradition
The Aboriginal people of the region around Ooldea have passed in oral tradition a detailed account of their mythology of the constellation Orion and the Pleiades. The anthropologist Daisy Bates was the first to attempt to compile records of this story. Work done by her and others has shown that all of the protagonists of the story of Nyeeruna and the Yugarilya correspond to individual stars covering the region around Orion and the Pleiades, with the exception of Baba, the father dingo, which is a major protagonist of the story and of the yearly re-enactments of the myth by the local people:
It has been suggested by Leaman and Hamacher that the location usually assigned to Baba by the locals (recorded by Bates as being at the "horn of the bull") is more likely to correspond to SN 1054 than to a faint star of that region such as β or ζ Tauri. This is motivated by the reference to Babba "returning to his place again" after attacking Nyeeruna which could refer to a transient star, as well as the fact that important characters of the myth are associated with bright stars. However, Leaman and Hamacher clarify there is no solid evidence to support this interpretation, which remains speculative. Hamacher demonstrates the extreme difficulty in identifying supernovae in Indigenous oral traditions.
Other elements of the story which have been found to correspond to astronomical elements by these authors include: awareness by the Aboriginal people of the different colors of the stars, possible awareness of the variability of Betelgeuse, observations of meteors in the Orionid meteor shower and the possibility that the rite associated with the myth is held at a time of astronomical significance, corresponding to the few days in the year when due to the Sun's proximity to Orion, it is unseen throughout the night, but is always in the sky during the daytime.
Media references
The supernova is mentioned in Ayreon's song To the Quasar, from the album Universal Migrator Part 2: Flight of the Migrator. SN 1054 and the lack of European recordings of the event is also mentioned in the historical fiction Space (Michener novel) by James A. Michener. The popular science book Death by Black Hole by Neil deGrasse Tyson uses SN 1054 to illustrate the relationships between religion, philosophy and human interpretations of astronomical events. The guest star of 1054 is also mentioned in Red Dragon (novel) by Thomas Harris.
| Physical sciences | Notable transient events | Astronomy |
1783069 | https://en.wikipedia.org/wiki/Time%20derivative | Time derivative | A time derivative is a derivative of a function with respect to time, usually interpreted as the rate of change of the value of the function. The variable denoting time is usually written as .
Notation
A variety of notations are used to denote the time derivative. In addition to the normal (Leibniz's) notation,
A very common short-hand notation used, especially in physics, is the 'over-dot'. I.E.
(This is called Newton's notation)
Higher time derivatives are also used: the second derivative with respect to time is written as
with the corresponding shorthand of .
As a generalization, the time derivative of a vector, say:
is defined as the vector whose components are the derivatives of the components of the original vector. That is,
Use in physics
Time derivatives are a key concept in physics. For example, for a changing position , its time derivative is its velocity, and its second derivative with respect to time, , is its acceleration. Even higher derivatives are sometimes also used: the third derivative of position with respect to time is known as the jerk. See motion graphs and derivatives.
A large number of fundamental equations in physics involve first or second time derivatives of quantities. Many other fundamental quantities in science are time derivatives of one another:
force is the time derivative of momentum
power is the time derivative of energy
electric current is the time derivative of electric charge
and so on.
A common occurrence in physics is the time derivative of a vector, such as velocity or displacement. In dealing with such a derivative, both magnitude and orientation may depend upon time.
Example: circular motion
For example, consider a particle moving in a circular path. Its position is given by the displacement vector , related to the angle, θ, and radial distance, r, as defined in the figure:
For this example, we assume that . Hence, the displacement (position) at any time t is given by
This form shows the motion described by r(t) is in a circle of radius r because the magnitude of r(t) is given by
using the trigonometric identity and where is the usual Euclidean dot product.
With this form for the displacement, the velocity now is found. The time derivative of the displacement vector is the velocity vector. In general, the derivative of a vector is a vector made up of components each of which is the derivative of the corresponding component of the original vector. Thus, in this case, the velocity vector is:
Thus the velocity of the particle is nonzero even though the magnitude of the position (that is, the radius of the path) is constant. The velocity is directed perpendicular to the displacement, as can be established using the dot product:
Acceleration is then the time-derivative of velocity:
The acceleration is directed inward, toward the axis of rotation. It points opposite to the position vector and perpendicular to the velocity vector. This inward-directed acceleration is called centripetal acceleration.
In differential geometry
In differential geometry, quantities are often expressed with respect to the local covariant basis, , where i ranges over the number of dimensions. The components of a vector expressed this way transform as a contravariant tensor, as shown in the expression , invoking Einstein summation convention. If we want to calculate the time derivatives of these components along a trajectory, so that we have , we can define a new operator, the invariant derivative , which will continue to return contravariant tensors:
where (with being the jth coordinate) captures the components of the velocity in the local covariant basis, and are the Christoffel symbols for the coordinate system. Note that explicit dependence on t has been repressed in the notation. We can then write:
as well as:
In terms of the covariant derivative, , we have:
Use in economics
In economics, many theoretical models of the evolution of various economic variables are constructed in continuous time and therefore employ time derivatives. One situation involves a stock variable and its time derivative, a flow variable. Examples include:
The flow of net fixed investment is the time derivative of the capital stock.
The flow of inventory investment is the time derivative of the stock of inventories.
The growth rate of the money supply is the time derivative of the money supply divided by the money supply itself.
Sometimes the time derivative of a flow variable can appear in a model:
The growth rate of output is the time derivative of the flow of output divided by output itself.
The growth rate of the labor force is the time derivative of the labor force divided by the labor force itself.
And sometimes there appears a time derivative of a variable which, unlike the examples above, is not measured in units of currency:
The time derivative of a key interest rate can appear.
The inflation rate is the growth rate of the price level—that is, the time derivative of the price level divided by the price level itself.
| Mathematics | Differential calculus | null |
1784313 | https://en.wikipedia.org/wiki/Tests%20of%20general%20relativity | Tests of general relativity | Tests of general relativity serve to establish observational evidence for the theory of general relativity. The first three tests, proposed by Albert Einstein in 1915, concerned the "anomalous" precession of the perihelion of Mercury, the bending of light in gravitational fields, and the gravitational redshift. The precession of Mercury was already known; experiments showing light bending in accordance with the predictions of general relativity were performed in 1919, with increasingly precise measurements made in subsequent tests; and scientists claimed to have measured the gravitational redshift in 1925, although measurements sensitive enough to actually confirm the theory were not made until 1954. A more accurate program starting in 1959 tested general relativity in the weak gravitational field limit, severely limiting possible deviations from the theory.
In the 1970s, scientists began to make additional tests, starting with Irwin Shapiro's measurement of the relativistic time delay in radar signal travel time near the Sun. Beginning in 1974, Hulse, Taylor and others studied the behaviour of binary pulsars experiencing much stronger gravitational fields than those found in the Solar System. Both in the weak field limit (as in the Solar System) and with the stronger fields present in systems of binary pulsars the predictions of general relativity have been extremely well tested.
In February 2016, the Advanced LIGO team announced that they had directly detected gravitational waves from a black hole merger. This discovery, along with additional detections announced in June 2016 and June 2017, tested general relativity in the very strong field limit, observing to date no deviations from theory.
Classical tests
Albert Einstein proposed three tests of general relativity, subsequently called the "classical tests" of general relativity, in 1916:
the perihelion precession of Mercury's orbit
the deflection of light by the Sun
the gravitational redshift of light
In the letter to The Times (of London) on November 28, 1919, he described the theory of relativity and thanked his English colleagues for their understanding and testing of his work. He also mentioned three classical tests with comments:
"The chief attraction of the theory lies in its logical completeness. If a single one of the conclusions drawn from it proves wrong, it must be given up; to modify it without destroying the whole structure seems to be impossible."
Perihelion precession of Mercury
Under Newtonian physics, an object in an (isolated) two-body system, consisting of the object orbiting a spherical mass, would trace out an ellipse with the center of mass of the system at a focus of the ellipse. The point of closest approach, called the periapsis (or when the central body is the Sun, perihelion), is fixed. Hence the major axis of the ellipse remains fixed in space. Both objects orbit around the center of mass of this system, so they each have their own ellipse. However, a number of effects in the Solar System cause the perihelia of planets to precess (rotate) around the Sun in the plane of their orbits, or equivalently, cause the major axis to rotate about the center of mass, hence changing its orientation in space. The principal cause is the presence of other planets which perturb one another's orbit. Another (much less significant) effect is solar oblateness.
Mercury deviates from the precession predicted from these Newtonian effects. This anomalous rate of precession of the perihelion of Mercury's orbit was first recognized in 1859 as a problem in celestial mechanics, by Urbain Le Verrier. His re-analysis of available timed observations of transits of Mercury over the Sun's disk from 1697 to 1848 showed that the actual rate of the precession disagreed from that predicted from Newton's theory by 38″ (arcseconds) per tropical century (later re-estimated at 43″ by Simon Newcomb in 1882). A number of ad hoc and ultimately unsuccessful solutions were proposed, but they tended to introduce more problems. Le Verrier suggested that another hypothetical planet might exist to account for Mercury's behavior. The previously successful search for Neptune based on its perturbations of the orbit of Uranus led astronomers to place some faith in this possible explanation, and the hypothetical planet was even named Vulcan. Finally, in 1908, W. W. Campbell, Director of the Lick Observatory, after the comprehensive photographic observations by Lick astronomer, Charles D. Perrine, at three solar eclipse expeditions, stated, "In my opinion, Dr. Perrine's work at the three eclipses of 1901, 1905, and 1908 brings the observational side of the famous intramercurial-planet problem definitely to a close." Subsequently, no evidence of Vulcan was found and Einstein's 1915 general theory accounted for Mercury's anomalous precession. Einstein wrote to Michele Besso, "Perihelion motions explained quantitatively ... you will be astonished".
In general relativity, this remaining precession, or change of orientation of the orbital ellipse within its orbital plane, is explained by gravitation being mediated by the curvature of spacetime. Einstein showed that general relativity agrees closely with the observed amount of perihelion shift. This was a powerful factor motivating the adoption of general relativity.
Although earlier measurements of planetary orbits were made using conventional telescopes, more accurate measurements are now made with radar. The total observed precession of Mercury is (574.10 ± 0.65)″ per century relative to the inertial ICRF. This precession can be attributed to the following causes:
The correction by ()″/cy is the prediction of post-Newtonian theory with parameters . Thus the effect can be fully explained by general relativity. More recent calculations based on more precise measurements have not materially changed the situation.
In general relativity the perihelion shift σ, expressed in radians per revolution, is approximately given by:
where L is the semi-major axis, T is the orbital period, c is the speed of light, and e is the orbital eccentricity (see: Two-body problem in general relativity).
The other planets experience perihelion shifts as well, but, since they are farther from the Sun and have longer periods, their shifts are lower, and could not be observed accurately until long after Mercury's. For example, the perihelion shift of Earth's orbit due to general relativity is theoretically 3.83868″ per century and experimentally ()″/cy, Venus's is 8.62473″/cy and (8.6247 ± 0.0005)″/cy and Mars' is ()″/cy. Both values have now been measured, with results in good agreement with theory. The periapsis shift has also now been measured for binary pulsar systems, with PSR 1913+16 amounting to 4.2° per year. These observations are consistent with general relativity. It is also possible to measure periapsis shift in binary star systems which do not contain ultra-dense stars, but it is more difficult to model the classical effects precisely – for example, the alignment of the stars' spin to their orbital plane needs to be known and is hard to measure directly. A few systems, such as DI Herculis, have been measured as test cases for general relativity.
Deflection of light by the Sun
Henry Cavendish in 1784 (in an unpublished manuscript) and Johann Georg von Soldner in 1801 (published in 1804) had pointed out that Newtonian gravity predicts that starlight will bend around a massive object. The same value as Soldner's was calculated by Einstein in 1911 based on the equivalence principle alone. However, Einstein noted in 1915 in the process of completing general relativity, that his 1911 result (and thus Soldner's 1801 result) is only half of the correct value. Einstein became the first to calculate the correct value for light bending: 1.75 arcseconds for light that grazes the Sun.
The first observation of light deflection was performed by noting the change in position of stars as they passed near the Sun on the celestial sphere. The observations were performed by Arthur Eddington and his collaborators (see Eddington experiment) during the total solar eclipse of May 29, 1919, when the stars near the Sun (at that time in the constellation Taurus) could be observed. Observations were made simultaneously in the cities of Sobral, Ceará, Brazil and in São Tomé and Príncipe on the west coast of Africa. The result was considered spectacular news and made the front page of most major newspapers. It made Einstein and his theory of general relativity world-famous. When asked by his assistant what his reaction would have been if general relativity had not been confirmed by Eddington and Dyson in 1919, Einstein famously made the quip: "Then I would feel sorry for the dear Lord. The theory is correct anyway."
The early accuracy, however, was poor and there was doubt that the small number of measured star locations and instrument questions could produce a reliable result. The results were argued by some to have been plagued by systematic error and possibly confirmation bias, although modern reanalysis of the dataset suggests that Eddington's analysis was accurate. The measurement was repeated by a team from the Lick Observatory led by the Director W. W. Campbell in the 1922 eclipse as observed in remote Australian station of Wallal, with results based on hundreds of star positions that agreed with the 1919 results and has been repeated several times since, most notably in 1953 by Yerkes Observatory astronomers and in 1973 by a team from the University of Texas. Considerable uncertainty remained in these measurements for almost fifty years, until observations started being made at radio frequencies.
The deflection of starlight by the nearby white dwarf star Stein 2051 B has also been measured.
Gravitational redshift of light
Einstein predicted the gravitational redshift of light from the equivalence principle in 1907, and it was predicted that this effect might be measured in the spectral lines of a white dwarf star, which has a very high gravitational field. Initial attempts to measure the gravitational redshift of the spectrum of Sirius-B, were done by Walter Sydney Adams in 1925, but the result was criticized as being unusable due to the contamination from light from the (much brighter) primary star, Sirius. The first accurate measurement of the gravitational redshift of a white dwarf was done by Popper in 1954, measuring a 21 km/s gravitational redshift of 40 Eridani B.
The redshift of Sirius B was finally measured by Greenstein et al. in 1971, obtaining the value for the gravitational redshift of , with more accurate measurements by the Hubble Space Telescope showing .
Tests of special relativity
The general theory of relativity incorporates Einstein's special theory of relativity, and hence tests of special relativity are also testing aspects of general relativity. As a consequence of the equivalence principle, Lorentz invariance holds locally in non-rotating, freely falling reference frames. Experiments related to Lorentz invariance special relativity (that is, when gravitational effects can be neglected) are described in tests of special relativity.
Modern tests
The modern era of testing general relativity was ushered in largely at the impetus of Dicke and Schiff who laid out a framework for testing general relativity. They emphasized the importance not only of the classical tests, but of null experiments, testing for effects which in principle could occur in a theory of gravitation, but do not occur in general relativity. Other important theoretical developments included the inception of alternative theories to general relativity, in particular, scalar–tensor theories such as the Brans–Dicke theory; the parameterized post-Newtonian formalism in which deviations from general relativity can be quantified; and the framework of the equivalence principle.
Experimentally, new developments in space exploration, electronics and condensed matter physics have made additional precise experiments possible, such as the Pound–Rebka experiment, laser interferometry and lunar rangefinding.
Post-Newtonian tests of gravity
Early tests of general relativity were hampered by the lack of viable competitors to the theory: it was not clear what sorts of tests would distinguish it from its competitors. General relativity was the only known relativistic theory of gravity compatible with special relativity and observations. Moreover, it is an extremely simple and elegant theory. This changed with the introduction of Brans–Dicke theory in 1960. This theory is arguably simpler, as it contains no dimensionful constants, and is compatible with a version of Mach's principle and Dirac's large numbers hypothesis, two philosophical ideas which have been influential in the history of relativity. Ultimately, this led to the development of the parametrized post-Newtonian formalism by Nordtvedt and Will, which parametrizes, in terms of ten adjustable parameters, all the possible departures from Newton's law of universal gravitation to first order in the velocity of moving objects (i.e. to first order in , where v is the velocity of an object and c is the speed of light). This approximation allows the possible deviations from general relativity, for slowly moving objects in weak gravitational fields, to be systematically analyzed. Much effort has been put into constraining the post-Newtonian parameters, and deviations from general relativity are at present severely limited.
The experiments testing gravitational lensing and light time delay limits the same post-Newtonian parameter, the so-called Eddington parameter γ, which is a straightforward parametrization of the amount of deflection of light by a gravitational source. It is equal to one for general relativity, and takes different values in other theories (such as Brans–Dicke theory). It is the best constrained of the ten post-Newtonian parameters, but there are other experiments designed to constrain the others. Precise observations of the perihelion shift of Mercury constrain other parameters, as do tests of the strong equivalence principle.
One of the goals of the BepiColombo mission to Mercury, is to test the general relativity theory by measuring the parameters gamma and beta of the parametrized post-Newtonian formalism with high accuracy. The experiment is part of the Mercury Orbiter Radio science Experiment (MORE). The spacecraft was launched in October 2018 and is expected to enter orbit around Mercury in December 2025.
Gravitational lensing
One of the most important tests is gravitational lensing. It has been observed in distant astrophysical sources, but these are poorly controlled and it is uncertain how they constrain general relativity. The most precise tests are analogous to Eddington's 1919 experiment: they measure the deflection of radiation from a distant source by the Sun. The sources that can be most precisely analyzed are distant radio sources. In particular, some quasars are very strong radio sources. The directional resolution of any telescope is in principle limited by diffraction; for radio telescopes this is also the practical limit. An important improvement in obtaining positional high accuracies (from milli-arcsecond to micro-arcsecond) was obtained by combining radio telescopes across Earth. The technique is called very long baseline interferometry (VLBI). With this technique radio observations couple the phase information of the radio signal observed in telescopes separated over large distances. Recently, these telescopes have measured the deflection of radio waves by the Sun to extremely high precision, confirming the amount of deflection predicted by general relativity aspect to the 0.03% level. At this level of precision systematic effects have to be carefully taken into account to determine the precise location of the telescopes on Earth. Some important effects are Earth's nutation, rotation, atmospheric refraction, tectonic displacement and tidal waves. Another important effect is refraction of the radio waves by the solar corona. Fortunately, this effect has a characteristic spectrum, whereas gravitational distortion is independent of wavelength. Thus, careful analysis, using measurements at several frequencies, can subtract this source of error.
The entire sky is slightly distorted due to the gravitational deflection of light caused by the Sun (the anti-Sun direction excepted). This effect has been observed by the European Space Agency astrometric satellite Hipparcos. It measured the positions of about 105 stars. During the full mission about relative positions have been determined, each to an accuracy of typically 3 milliarcseconds (the accuracy for an 8–9 magnitude star). Since the gravitation deflection perpendicular to the Earth–Sun direction is already 4.07 milliarcseconds, corrections are needed for practically all stars. Without systematic effects, the error in an individual observation of 3 milliarcseconds, could be reduced by the square root of the number of positions, leading to a precision of 0.0016 milliarcseconds. Systematic effects, however, limit the accuracy of the determination to 0.3% (Froeschlé, 1997).
Launched in 2013, the Gaia spacecraft will conduct a census of one billion stars in the Milky Way and measure their positions to an accuracy of 24 microarcseconds. Thus it will also provide stringent new tests of gravitational deflection of light caused by the Sun which was predicted by General relativity.
Light travel time delay testing
Irwin I. Shapiro proposed another test, beyond the classical tests, which could be performed within the Solar System. It is sometimes called the fourth "classical" test of general relativity. He predicted a relativistic time delay (Shapiro delay) in the round-trip travel time for radar signals reflecting off other planets. The mere curvature of the path of a photon passing near the Sun is too small to have an observable delaying effect (when the round-trip time is compared to the time taken if the photon had followed a straight path), but general relativity predicts a time delay that becomes progressively larger when the photon passes nearer to the Sun due to the time dilation in the gravitational potential of the Sun. Observing radar reflections from Mercury and Venus just before and after they are eclipsed by the Sun agrees with general relativity theory at the 5% level.
More recently, the Cassini probe has undertaken a similar experiment which gave agreement with general relativity at the 0.002% level. However, the following detailed studies revealed that the measured value of the PPN parameter gamma is affected by a gravitomagnetic effect caused by the orbital motion of Sun around the barycenter of the solar system. The gravitomagnetic effect in the Cassini radioscience experiment was implicitly postulated by B. Bertotti as having a pure general relativistic origin but its theoretical value has never been tested in the experiment which effectively makes the experimental uncertainty in the measured value of gamma actually larger (by a factor of 10) than 0.002% claimed by B. Bertotti and co-authors in Nature.
Very Long Baseline Interferometry has measured velocity-dependent (gravitomagnetic) corrections to the Shapiro time delay in the field of moving Jupiter and Saturn.
Equivalence principle
The equivalence principle, in its simplest form, asserts that the trajectories of falling bodies in a gravitational field should be independent of their mass and internal structure, provided they are small enough not to disturb the environment or be affected by tidal forces. This idea has been tested to extremely high precision by Eötvös torsion balance experiments, which look for a differential acceleration between two test masses. Constraints on this, and on the existence of a composition-dependent fifth force or gravitational Yukawa interaction are very strong, and are discussed under fifth force and weak equivalence principle.
A version of the equivalence principle, called the strong equivalence principle, asserts that self-gravitation falling bodies, such as stars, planets or black holes (which are all held together by their gravitational attraction) should follow the same trajectories in a gravitational field, provided the same conditions are satisfied. This is called the Nordtvedt effect and is most precisely tested by the Lunar Laser Ranging Experiment. Since 1969, it has continuously measured the distance from several rangefinding stations on Earth to reflectors on the Moon to approximately centimeter accuracy. These have provided a strong constraint on several of the other post-Newtonian parameters.
Another part of the strong equivalence principle is the requirement that Newton's gravitational constant be constant in time, and have the same value everywhere in the universe. There are many independent observations limiting the possible variation of Newton's gravitational constant, but one of the best comes from lunar rangefinding which suggests that the gravitational constant does not change by more than one part in 1011 per year. The constancy of the other constants is discussed in the Einstein equivalence principle section of the equivalence principle article.
Gravitational redshift and time dilation
The first of the classical tests discussed above, the gravitational redshift, is a simple consequence of the Einstein equivalence principle and was predicted by Einstein in 1907. As such, it is not a test of general relativity in the same way as the post-Newtonian tests, because any theory of gravity obeying the equivalence principle should also incorporate the gravitational redshift. Nonetheless, confirming the existence of the effect was an important substantiation of relativistic gravity, since the absence of gravitational redshift would have strongly contradicted relativity. The first observation of the gravitational redshift was the measurement of the shift in the spectral lines from the white dwarf star Sirius B by Adams in 1925, discussed above, and follow-on measurements of other white dwarfs. Because of the difficulty of the astrophysical measurement, however, experimental verification using a known terrestrial source was preferable.
Experimental verification of gravitational redshift using terrestrial sources took several decades, because it is difficult to find clocks (to measure time dilation) or sources of electromagnetic radiation (to measure redshift) with a frequency that is known well enough that the effect can be accurately measured. It was confirmed experimentally for the first time in 1959 using measurements of the change in wavelength of gamma-ray photons generated with the Mössbauer effect, which generates radiation with a very narrow line width. The Pound–Rebka experiment measured the relative redshift of two sources situated at the top and bottom of Harvard University's Jefferson tower. The result was in excellent agreement with general relativity. This was one of the first precision experiments testing general relativity. The experiment was later improved to better than the 1% level by Pound and Snider.
The blueshift of a falling photon can be found by assuming it has an equivalent mass based on its frequency (where h is the Planck constant) along with , a result of special relativity. Such simple derivations ignore the fact that in general relativity the experiment compares clock rates, rather than energies. In other words, the "higher energy" of the photon after it falls can be equivalently ascribed to the slower running of clocks deeper in the gravitational potential well. To fully validate general relativity, it is important to also show that the rate of arrival of the photons is greater than the rate at which they are emitted. A very accurate gravitational redshift experiment, which deals with this issue, was performed in 1976, where a hydrogen maser clock on a rocket was launched to a height of 10,000 km, and its rate compared with an identical clock on the ground. It tested the gravitational redshift to 0.007%.
Although the Global Positioning System (GPS) is not designed as a test of fundamental physics, it must account for the gravitational redshift in its timing system, and physicists have analyzed timing data from the GPS to confirm other tests. When the first satellite was launched, some engineers resisted the prediction that a noticeable gravitational time dilation would occur, so the first satellite was launched without the clock adjustment that was later built into subsequent satellites. It showed the predicted shift of 38 microseconds per day. This rate of discrepancy is sufficient to substantially impair function of GPS within hours if not accounted for. An excellent account of the role played by general relativity in the design of GPS can be found in Ashby 2003.
Other precision tests of general relativity, not discussed here, are the Gravity Probe A satellite, launched in 1976, which showed gravity and velocity affect the ability to synchronize the rates of clocks orbiting a central mass and the Hafele–Keating experiment, which used atomic clocks in circumnavigating aircraft to test general relativity and special relativity together.
Frame-dragging tests
Tests of the Lense–Thirring precession, consisting of small secular precessions of the orbit of a test particle in motion around a central rotating mass, for example, a planet or a star, have been performed with the LAGEOS satellites, but many aspects of them remain controversial. The same effect may have been detected in the data of the Mars Global Surveyor (MGS) spacecraft, a former probe in orbit around Mars; also such a test raised a debate. First attempts to detect the Sun's Lense–Thirring effect on the perihelia of the inner planets have been recently reported as well. Frame dragging would cause the orbital plane of stars orbiting near a supermassive black hole to precess about the black hole spin axis. This effect should be detectable within the next few years via astrometric monitoring of stars at the center of the Milky Way galaxy. By comparing the rate of orbital precession of two stars on different orbits, it is possible in principle to test the no-hair theorems of general relativity.
The Gravity Probe B satellite, launched in 2004 and operated until 2005, detected frame-dragging and the geodetic effect. The experiment used four quartz spheres the size of ping pong balls coated with a superconductor. Data analysis continued through 2011 due to high noise levels and difficulties in modelling the noise accurately so that a useful signal could be found. Principal investigators at Stanford University reported on May 4, 2011, that they had accurately measured the frame dragging effect relative to the distant star IM Pegasi, and the calculations proved to be in line with the prediction of Einstein's theory. The results, published in Physical Review Letters measured the geodetic effect with an error of about 0.2 percent. The results reported the frame dragging effect (caused by Earth's rotation) added up to 37 milliarcseconds with an error of about 19 percent. Investigator Francis Everitt explained that a milliarcsecond "is the width of a human hair seen at the distance of 10 miles".
In January 2012, LARES satellite was launched on a Vega rocket to measure Lense–Thirring effect with an accuracy of about 1%, according to its proponents. This evaluation of the actual accuracy obtainable is a subject of debate.
Tests of the gravitational potential at small distances
It is possible to test whether the gravitational potential continues with the inverse square law at very small distances. Tests so far have focused on a divergence from GR in the form of a Yukawa potential , but no evidence for a potential of this kind has been found. The Yukawa potential with has been ruled out down to .
Mössbauer rotor experiment
It was conceived as a means to measure the time dilation effect on Earth after being motivated by Einstein's equivalence principle that implies a rotating observer will be subject to the same transformations as an observer in a gravitational field. Mössbauer rotor experiments hence permit a precise terrestrial test of the relativistic Doppler effect. From a radioactive source fixed at the center of a spinning disc or rod, gamma rays travel to an absorber at the rim (in some variations of the experiment this scheme was reversed) and an unabsorbed number of them pass through depending on the rotational speed to arrive at a stationary counter (i.e., detector of gamma quanta resting in the lab frame). In lieu with the Clock hypothesis, Einstein's general relativity predicts that the moving absorber's clock at the rim should retard by a specific amount due to time dilation on account of centrifugal binding alone compared to a rest frame absorber. So the transmission of gamma photons through the absorber should increase during rotation, which can be subsequently measured by the stationary counter beyond the absorber. This prediction was actually observed using the Mössbauer effect, since the equivalence principle, as originally suggested by Einstein, implicitly allows the association of the time dilation due to rotation (calculated as a result of the change in the detector's count rate) with gravitational time dilation. Such experiments were pioneered by Hay et al. (1960), Champeney et al. (1965), and Kündig (1963), and all of them had declared confirmation of the prediction of Einstein's theory of relativity.
Be that as it may, an early 21st Century re-examination of these endeavors called into question the validity of the past obtained results claiming to have verified time dilation as predicted by Einstein's relativity theory, whereby novel experimentations were carried out that uncovered an extra energy shift between emitted and absorbed radiation next to the classical relativistic dilation of time. This discovery was first explained as discrediting general relativity and successfully confirming at the laboratory scale the predictions of an alternative theory of gravity developed by T. Yarman and his colleagues. Against this development, a contentious attempt was made to explain the disclosed extra energy shift as arising from a so-far unknown and allegedly missed clock synchronization effect, which was unusually awarded a prize in 2018 by the Gravity Research Foundation for having secured a new proof of general relativity. However, at the same time period, it was revealed that said author committed several mathematical errors in his calculations, and the supposed contribution of the so-called clock synchronization to the measured time dilation is in fact practically null. As a consequence, a general relativistic explanation for the outcomes of Mössbauer rotor experiments remains open.
Strong field tests
The very strong gravitational fields that are present close to black holes, especially those supermassive black holes which are thought to power active galactic nuclei and the more active quasars, belong to a field of intense active research. Observations of these quasars and active galactic nuclei are difficult, and interpretation of the observations is heavily dependent upon astrophysical models other than general relativity or competing fundamental theories of gravitation, but they are qualitatively consistent with the black hole concept as modeled in general relativity.
Binary pulsars
Pulsars are rapidly rotating neutron stars which emit regular radio pulses as they rotate. As such they act as clocks which allow very precise monitoring of their orbital motions. Observations of pulsars in orbit around other stars have all demonstrated substantial periapsis precessions that cannot be accounted for classically but can be accounted for by using general relativity. For example, the Hulse–Taylor binary pulsar PSR B1913+16 (a pair of neutron stars in which one is detected as a pulsar) has an observed precession of over 4° of arc per year (periastron shift per orbit only about 10−6). This precession has been used to compute the masses of the components.
Similarly to the way in which atoms and molecules emit electromagnetic radiation, a gravitating mass that is in quadrupole type or higher order vibration, or is asymmetric and in rotation, can emit gravitational waves. These gravitational waves are predicted to travel at the speed of light. For example, planets orbiting the Sun constantly lose energy via gravitational radiation, but this effect is so small that it is unlikely it will be observed in the near future (Earth radiates about 200 watts of gravitational radiation).
The radiation of gravitational waves has been inferred from the Hulse–Taylor binary (and other binary pulsars). Precise timing of the pulses shows that the stars orbit only approximately according to Kepler's Laws: over time they gradually spiral towards each other, demonstrating an energy loss in close agreement with the predicted energy radiated by gravitational waves. For their discovery of the first binary pulsar and measuring its orbital decay due to gravitational-wave emission, Hulse and Taylor won the 1993 Nobel Prize in Physics.
A "double pulsar" discovered in 2003, PSR J0737-3039, has a periastron precession of 16.90° per year; unlike the Hulse–Taylor binary, both neutron stars are detected as pulsars, allowing precision timing of both members of the system. Due to this, the tight orbit, the fact that the system is almost edge-on, and the very low transverse velocity of the system as seen from Earth, J0737−3039 provides by far the best system for strong-field tests of general relativity known so far. Several distinct relativistic effects are observed, including orbital decay as in the Hulse–Taylor system. After observing the system for two and a half years, four independent tests of general relativity were possible, the most precise (the Shapiro delay) confirming the general relativity prediction within 0.05% (nevertheless the periastron shift per orbit is only about 0.0013% of a circle and thus it is not a higher-order relativity test).
In 2013, an international team of astronomers reported new data from observing a pulsar-white dwarf system PSR J0348+0432, in which they have been able to measure a change in the orbital period of 8 millionths of a second per year, and confirmed GR predictions in a regime of extreme gravitational fields never probed before; but there are still some competing theories that would agree with these data.
Direct detection of gravitational waves
A number of gravitational-wave detectors have been built with the intent of directly detecting the gravitational waves emanating from such astronomical events as the merger of two neutron stars or black holes. In February 2016, the Advanced LIGO team announced that they had directly detected gravitational waves from a stellar binary black hole merger, with additional detections announced in June 2016, June 2017, and August 2017.
General relativity predicts gravitational waves, as does any theory of gravitation in which changes in the gravitational field propagate at a finite speed. Then, the LIGO response function could discriminate among the various theories. Since gravitational waves can be directly detected, it is possible to use them to learn about the Universe. This is gravitational-wave astronomy. Gravitational-wave astronomy can test general relativity by verifying that the observed waves are of the form predicted (for example, that they only have two transverse polarizations), and by checking that black holes are the objects described by solutions of the Einstein field equations.
Gravitational-wave astronomy can also test Maxwell-Einstein field equations. This version of the field equations predicts that spinning magnetars (i.e., neutron stars with extremely strong magnetic dipole field) should emit gravitational waves.
"These amazing observations are the confirmation of a lot of theoretical work, including Einstein's general theory of relativity, which predicts gravitational waves", said Stephen Hawking.
Direct observation of black holes
The galaxy M87 was the subject of observation by the Event Horizon Telescope (EHT) in 2017; the 10 April 2019 issue of Astrophysical Journal Letters (vol. 875, No. 1) was dedicated to the EHT results, publishing six open-access papers. The event horizon of the black hole at the center of M87 was directly imaged at the wavelength of radio waves by the EHT; the image was revealed in a press conference on 10 April 2019, the first image of a black hole's event horizon.
In May 2022, the EHT provided the first image of the super massive black hole Sagittarius A* in the center of our own Milky Way galaxy.
Gravitational redshift and orbit precession of star in strong gravity field
Gravitational redshift in light from the S2 star orbiting the supermassive black hole Sagittarius A* in the center of the Milky Way has been measured with the Very Large Telescope using GRAVITY, NACO and SIFONI instruments. Additionally, there has now been detection of the Schwarzschild precession in the orbit of the star S2 near the Galactic centre massive black hole.
Strong equivalence principle
The strong equivalence principle of general relativity requires universality of free fall to apply even to bodies with strong self-gravity. Direct tests of this principle using Solar System bodies are limited by the weak self-gravity of the bodies, and tests using pulsar–white-dwarf binaries have been limited by the weak gravitational pull of the Milky Way. With the discovery of a triple star system called PSR J0337+1715, located about 4,200 light-years from Earth, the strong equivalence principle can be tested with a high accuracy. This system contains a neutron star in a 1.6-day orbit with a white dwarf star, and the pair in a 327-day orbit with another white dwarf further away. This system permits a test that compares how the gravitational pull of the outer white dwarf affects the pulsar, which has strong self-gravity, and the inner white dwarf. The result shows that the accelerations of the pulsar and its nearby white-dwarf companion differ fractionally by no more than 2.6 (95% confidence level).
X-ray spectroscopy
This technique is based on the idea that photon trajectories are modified in the presence of a gravitational body. A very common astrophysical system in the universe is a black hole surrounded by an accretion disk. The radiation from the general neighborhood, including the accretion disk, is affected by the nature of the central black hole. Assuming Einstein's theory is correct, astrophysical black holes are described by the Kerr metric. (A consequence of the no-hair theorems.) Thus, by analyzing the radiation from such systems, it is possible to test Einstein's theory.
Most of the radiation from these black hole – accretion disk systems (e.g., black hole binaries and active galactic nuclei) arrives in the form of X-rays. When modeled, the radiation is decomposed into several components. Tests of Einstein's theory are possible with the thermal spectrum (only for black hole binaries) and the reflection spectrum (for both black hole binaries and active galactic nuclei). The former is not expected to provide strong constraints, while the latter is much more promising. In both cases, systematic uncertainties might make such tests more challenging.
Cosmological tests
Tests of general relativity on the largest scales are not nearly so stringent as Solar System tests. The earliest such test was the prediction and discovery of the expansion of the universe. In 1922, Alexander Friedmann found that the Einstein equations have non-stationary solutions (even in the presence of the cosmological constant). In 1927, Georges Lemaître showed that static solutions of the Einstein equations, which are possible in the presence of the cosmological constant, are unstable, and therefore the static universe envisioned by Einstein could not exist (it must either expand or contract). Lemaître made an explicit prediction that the universe should expand. He also derived a redshift-distance relationship, which is now known as the Hubble Law. Later, in 1931, Einstein himself agreed with the results of Friedmann and Lemaître. The expansion of the universe discovered by Edwin Hubble in 1929 was then considered by many (and continues to be considered by some now) as a direct confirmation of general relativity. In the 1930s, largely due to the work of E. A. Milne, it was realised that the linear relationship between redshift and distance derives from the general assumption of uniformity and isotropy rather than specifically from general relativity. However the prediction of a non-static universe was non-trivial, indeed dramatic, and primarily motivated by general relativity.
Some other cosmological tests include searches for primordial gravitational waves generated during cosmic inflation, which may be detected in the cosmic microwave background polarization or by a proposed space-based gravitational-wave interferometer called the Big Bang Observer. Other tests at high redshift are constraints on other theories of gravity, and the variation of the gravitational constant since Big Bang nucleosynthesis (it varied by no more than 40% since then).
In August 2017, the findings of tests conducted by astronomers using the European Southern Observatory's Very Large Telescope (VLT), among other instruments, were released, and positively demonstrated gravitational effects predicted by Albert Einstein. One of these tests observed the orbit of the stars circling around Sagittarius A*, a black hole about 4 million times as massive as the sun. Einstein's theory suggested that large objects bend the space around them, causing other objects to diverge from the straight lines they would otherwise follow. Although previous studies have validated Einstein's theory, this was the first time his theory had been tested on such a gigantic object. The findings were published in The Astrophysical Journal.
Gravitational lensing
Astronomers using the Hubble Space Telescope and the Very Large Telescope have made precise tests of general relativity on galactic scales. The nearby galaxy ESO 325-G004 acts as a strong gravitational lens, distorting light from a distant galaxy behind it to create an Einstein ring around its centre. By comparing the mass of ESO 325-G004 (from measurements of the motions of stars inside this galaxy) with the curvature of space around it, astronomers found that gravity behaves as predicted by general relativity on these astronomical length-scales.
| Physical sciences | Theory of relativity | Physics |
1785475 | https://en.wikipedia.org/wiki/Chain-growth%20polymerization | Chain-growth polymerization | Chain-growth polymerization (AE) or chain-growth polymerisation (BE) is a polymerization technique where monomer molecules add onto the active site on a growing polymer chain one at a time. There are a limited number of these active sites at any moment during the polymerization which gives this method its key characteristics.
Chain-growth polymerization involves 3 types of reactions :
Initiation: An active species I* is formed by some decomposition of an initiator molecule I
Propagation: The initiator fragment reacts with a monomer M to begin the conversion to the polymer; the center of activity is retained in the adduct. Monomers continue to add in the same way until polymers Pi* are formed with the degree of polymerization i
Termination: By some reaction generally involving two polymers containing active centers, the growth center is deactivated, resulting in dead polymer
Introduction
In 1953, Paul Flory first classified polymerization as "step-growth polymerization" and "chain-growth polymerization". IUPAC recommends to further simplify "chain-growth polymerization" to "chain polymerization". It is a kind of polymerization where an active center (free radical or ion) is formed, and a plurality of monomers can be polymerized together in a short period of time to form a macromolecule having a large molecular weight. In addition to the regenerated active sites of each monomer unit, polymer growth will only occur at one (or possibly more) endpoint.
Many common polymers can be obtained by chain polymerization such as polyethylene (PE), polypropylene (PP), polyvinyl chloride (PVC), poly(methyl methacrylate) (PMMA), polyacrylonitrile (PAN), polyvinyl acetate (PVA).
Typically, chain-growth polymerization can be understood with the chemical equation:
In this equation, P is the polymer while x represents degree of polymerization, * means active center of chain-growth polymerization, M is the monomer which will react with active center, and L may be a low-molar-mass by-product obtained during chain propagation. For most chain-growth polymerizations, there is no by-product L formed. However there are some exceptions, such as the polymerization of amino acid N-carboxyanhydrides to oxazolidine-2,5-diones.
This type of polymerization is described as "chain" or "chain-growth" because the reaction mechanism is a chemical chain reaction with an initiation step in which an active center is formed, followed by a rapid sequence of chain propagation steps in which the polymer molecule grows by addition of one monomer molecule to the active center in each step. The word "chain" here does not refer to the fact that polymer molecules form long chains. Some polymers are formed instead by a second type of mechanism known as step-growth polymerization without rapid chain propagation steps.
Reaction steps
All chain-growth polymerization reactions must include chain initiation and chain propagation. Chain transfer and chain termination steps also occur in many but not all chain-growth polymerizations.
Chain initiation
Chain initiation is the initial generation of a chain carrier, which is an intermediate such as a radical or an ion which can continue the reaction by chain propagation. Initiation steps are classified according to the way that energy is provided: thermal initiation, high energy initiation, and chemical initiation, etc. Thermal initiation uses molecular thermal motion to dissociate a molecule and form active centers. High energy initiation refers to the generation of chain carriers by radiation. Chemical initiation is due to a chemical initiator.
For the case of radical polymerization as an example, chain initiation involves the dissociation of a radical initiator molecule (I) which is easily dissociated by heat or light into two free radicals (2 R°). Each radical R° then adds a first monomer molecule (M) to start a chain which terminates with a monomer activated by the presence of an unpaired electron (RM1°).
I → 2 R°
R° + M → RM1°
Chain propagation
IUPAC defines chain propagation as a reaction of an active center on the growing polymer molecule, which adds one monomer molecule to form a new polymer molecule (RM1°) one repeat unit longer.
For radical polymerization, the active center remains an atom with an unpaired electron. The addition of the second monomer and a typical later addition step are
RM1° + M → RM2°
...............
RMn° + M → RMn+1°
For some polymers, chains of over 1000 monomer units can be formed in milliseconds.
Chain termination
In a chain termination step, the active center disappears, resulting in the termination of chain propagation. This is different from chain transfer in which the active center only shifts to another molecule but does not disappear.
For radical polymerization, termination involves a reaction of two growing polymer chains to eliminate the unpaired electrons of both chains. There are two possibilities.
1. Recombination is the reaction of the unpaired electrons of two chains to form a covalent bond between them. The product is a single polymer molecule with the combined length of the two reactant chains:
RMn° + RMm° → Pn+m
2. Disproportionation is the transfer of a hydrogen atom from one chain to the other, so that the two product chain molecules are unchanged in length but are no longer free radicals:
RMn° + RMm° → Pn + Pm
Initiation, propagation and termination steps also occur in chain reactions of smaller molecules. This is not true of the chain transfer and branching steps considered next.
Chain transfer
In some chain-growth polymerizations there is also a chain transfer step, in which the growing polymer chain RMn° takes an atom X from an inactive molecule XY, terminating the growth of the polymer chain: RMn° + XY → RMnX + Y°. The Y fragment ls a new active center which adds more monomer M to form a new growing chain YMn°. This can happen in free radical polymerization for chains RMn°, in ionic polymerization for chains RMn+ or RMn–, or in coordination polymerization. In most cases chain transfer will generate a by-product and decrease the molar mass of the final polymer.
Chain transfer to polymer: Branching
Another possibility is chain transfer to a second polymer molecule, result in the formation of a product macromolecule with a branched structure. In this case the growing chain takes an atom X from a second polymer chain whose growth had been completed. The growth of the first polymer chain is completed by the transfer of atom X. However the second molecule loses an atom X from the interior of its polymer chain to form a reactive radical (or ion) which can add more monomer molecules. This results in the addition of a branch or side chain and the formation of a product macromolecule with a branched structure.
Classes of chain-growth polymerization
The International Union of Pure and Applied Chemistry (IUPAC) recommends definitions for several classes of chain-growth polymerization.
Radical polymerization
Based on the IUPAC definition, radical polymerization is a chain polymerization in which the kinetic-chain carriers are radicals. Usually, the growing chain end bears an unpaired electron. Free radicals can be initiated by many methods such as heating, redox reactions, ultraviolet radiation, high energy irradiation, electrolysis, sonication, and plasma.
Free radical polymerization is very important in polymer chemistry. It is one of the most developed methods in chain-growth polymerization. Currently, most polymers in our daily life are synthesized by free radical polymerization, including polyethylene, polystyrene, polyvinyl chloride, polymethyl methacrylate, polyacrylonitrile, polyvinyl acetate, styrene butadiene rubber, nitrile rubber, neoprene, etc.
Ionic polymerization
Ionic polymerization is a chain polymerization in which the kinetic-chain carriers are ions or ion pairs. It can be further divided into anionic polymerization and cationic polymerization.
Ionic polymerization generates many polymers used in daily life, such as butyl rubber, polyisobutylene, polyphenylene, polyoxymethylene, polysiloxane, polyethylene oxide, high density polyethylene, isotactic polypropylene, butadiene rubber, etc. Living anionic polymerization was developed in the 1950s. The chain will remain active indefinitely unless the reaction is transferred or terminated deliberately, which allows the control of molar weight and dispersity (or polydispersity index, PDI).
Coordination polymerization
Coordination polymerization is a chain polymerization that involves the preliminary coordination of a monomer molecule with a chain carrier. The monomer is first coordinated with the transition metal active center, and then the activated monomer is inserted into the transition metal-carbon bond for chain growth. In some cases, coordination polymerization is also called insertion polymerization or complexing polymerization.
Advanced coordination polymerizations can control the tacticity, molecular weight and PDI of the polymer effectively. In addition, the racemic mixture of the chiral metallocene can be separated into its enantiomers. The oligomerization reaction produces an optically active branched olefin using an optically active catalyst.
Living polymerization
Living polymerization was first described by Michael Szwarc in 1956. It is defined as a chain polymerization from which chain transfer and chain termination are absent. In the absence of chain-transfer and chain termination, the monomer in the system is consumed and the polymerization stops but the polymer chain remains active. If new monomer is added, the polymerization can proceed.
Due to the low PDI and predictable molecular weight, living polymerization is at the forefront of polymer research. It can be further divided into living free radical polymerization, living ionic polymerization and living ring-opening metathesis polymerization, etc.
Ring-opening polymerization
Ring-opening polymerization is defined as a polymerization in which a cyclic monomer yields a monomeric unit which is acyclic or contains fewer cycles than the monomer. Generally, the ring-opening polymerization is carried out under mild conditions, and the by-product is less than in the polycondensation reaction. A high molecular weight polymer is easily obtained.
Common ring-opening polymerization products includes polypropylene oxide, polytetrahydrofuran, , polyoxymethylene, polycaprolactam and polysiloxane.
Reversible-deactivation polymerization
Reversible-deactivation polymerization is defined as a chain polymerization propagated by chain carriers that are deactivated reversibly, bringing them into one or more active-dormant equilibria. An example of a reversible-deactivation polymerization is group-transfer polymerization.
Comparison with step-growth polymerization
Polymers were first classified according to polymerization method by Wallace Carothers in 1929, who introduced the terms addition polymer and condensation polymer to describe polymers made by addition reactions and condensation reactions respectively. However this classification is inadequate to describe a polymer which can be made by either type of reaction, for example nylon 6 which can be made either by addition of a cyclic monomer or by condensation of a linear monomer.
Flory revised the classification to chain-growth polymerization and step-growth polymerization, based on polymerization mechanisms rather than polymer structures. IUPAC now recommends that the names of step-growth polymerization and chain-growth polymerization be further simplified to polycondensation (or polyaddition if no low-molar-mass by-product is formed when a monomer is added) and chain polymerization.
Most polymerizations are either chain-growth or step-growth reactions. Chain-growth includes both initiation and propagation steps (at least), and the propagation of chain-growth polymers proceeds by the addition of monomers to a growing polymer with an active centre. In contrast step-growth polymerization involves only one type of step, and macromolecules can grow by reaction steps between any two molecular species: two monomers, a monomer and a growing chain, or two growing chains. In step growth, the monomers will initially form dimers, trimers, etc. which later react to form long chain polymers.
In chain-growth polymerization, a growing macromolecule increases in size rapidly once its growth is initiated. When a macromolecule stops growing it generally will add no more monomers. In step-growth polymerization on the other hand, a single polymer molecule can grow over the course of the whole reaction.
In chain-growth polymerization, long macromolecules with high molecular weight are formed when only a small fraction of monomer has reacted. Monomers are consumed steadily over the course of the whole reaction, but the degree of polymerization can increase very quickly after chain initiation. However in step-growth polymerization the monomer is consumed very quickly to dimer, trimer and oligomer. The degree of polymerization increases steadily during the whole polymerization process.
The type of polymerization of a given monomer usually depends on the functional groups present, and sometimes also on whether the monomer is linear or cyclic. Chain-growth polymers are usually addition polymers by Carothers' definition. They are typically formed by addition reactions of C=C bonds in the monomer backbone, which contains only carbon-carbon bonds. Another possibility is ring-opening polymerization, as for the chain-growth polymerization of tetrahydrofuran or of polycaprolactone (see Introduction above).
Step-growth polymers are typically condensation polymers in which an elimination product as such as H2O are formed. Examples are polyamides, polycarbonates, polyesters, polyimides, polysiloxanes and polysulfones. If no elimination product is formed, then the polymer is an addition polymer, such as a polyurethane or a poly(phenylene oxide). Chain-growth polymerization with a low-molar-mass by-product during chain growth is described by IUPAC as "condensative chain polymerization".
Compared to step-growth polymerization, living chain-growth polymerization shows low molar mass dispersity (or PDI), predictable molar mass distribution and controllable conformation. Generally, polycondensation proceeds in a step-growth polymerization mode.
Application
Chain polymerization products are widely used in many aspects of life, including electronic devices, food packaging, catalyst carriers, medical materials, etc. At present, the world's highest yielding polymers such as polyethylene (PE), polyvinyl chloride (PVC), polypropylene (PP), etc. can be obtained by chain polymerization.
In addition, some carbon nanotube polymer is used for electronical devices. Controlled living chain-growth conjugated polymerization will also enable the synthesis of well-defined advanced structures, including block copolymers. Their industrial applications extend to water purification, biomedical devices and sensors.
| Physical sciences | Organic reactions | Chemistry |
1786719 | https://en.wikipedia.org/wiki/Step-growth%20polymerization | Step-growth polymerization | In polymer chemistry, step-growth polymerization refers to a type of polymerization mechanism in which bi-functional or multifunctional monomers react to form first dimers, then trimers, longer oligomers and eventually long chain polymers. Many naturally-occurring and some synthetic polymers are produced by step-growth polymerization, e.g. polyesters, polyamides, polyurethanes, etc. Due to the nature of the polymerization mechanism, a high extent of reaction is required to achieve high molecular weight. The easiest way to visualize the mechanism of a step-growth polymerization is a group of people reaching out to hold their hands to form a human chain—each person has two hands (= reactive sites). There also is the possibility to have more than two reactive sites on a monomer: In this case branched polymers production take place.
IUPAC has deprecated the term step-growth polymerization, and recommends use of the terms polyaddition (when the propagation steps are addition reactions and molecules are not evolved during these steps) and polycondensation (when the propagation steps are condensation reactions and molecules are evolved during these steps).
Historical aspects
Most natural polymers being employed at early stage of human society are of condensation type. The synthesis of first truly synthetic polymeric material, bakelite, was announced by Leo Baekeland in 1907, through a typical step-growth polymerization fashion of phenol and formaldehyde.
The pioneer of synthetic polymer science, Wallace Carothers, developed a new means of making polyesters through step-growth polymerization in 1930s as a research group leader at DuPont. It was the first reaction designed and carried out with the specific purpose of creating high molecular weight polymer molecules, as well as the first polymerization reaction whose results had been predicted by scientific theory. Carothers developed a series of mathematic equations to describe the behavior of step-growth polymerization systems which are still known as the Carothers equations today. Collaborating with Paul Flory, a physical chemist, they developed theories that describe more mathematical aspects of step-growth polymerization including kinetics, stoichiometry, and molecular weight distribution etc. Carothers is also well known for his invention of Nylon.
Condensation polymerization
"Step growth polymerization" and condensation polymerization are two different concepts, not always identical. In fact polyurethane polymerizes with addition polymerization (because its polymerization produces no small molecules), but its reaction mechanism corresponds to a step-growth polymerization.
The distinction between "addition polymerization" and "condensation polymerization" was introduced by Wallace Carothers in 1929, and refers to the type of products, respectively:
a polymer only (addition)
a polymer and a molecule with a low molecular weight (condensation)
The distinction between "step-growth polymerization" and "chain-growth polymerization" was introduced by Paul Flory in 1953, and refers to the reaction mechanisms, respectively:
by functional groups (step-growth polymerization)
by free-radical or ion (chain-growth polymerization)
Differences from chain-growth polymerization
This technique is usually compared with chain-growth polymerization to show its characteristics.
Classes of step-growth polymers
Classes of step-growth polymers are:
Polyester has high glass transition temperature Tg and high melting point Tm, good mechanical properties to about 175 °C, good resistance to solvent and chemicals. It can exist as fibers and films. The former is used in garments, felts, tire cords, etc. The latter appears in magnetic recording tape and high grade films.
Polyamide (nylon) has good balance of properties: high strength, good elasticity and abrasion resistance, good toughness, favorable solvent resistance. The applications of polyamide include: rope, belting, fiber cloths, thread, substitute for metal in bearings, jackets on electrical wire.
Polyurethane can exist as elastomers with good abrasion resistance, hardness, good resistance to grease and good elasticity, as fibers with excellent rebound, as coatings with good resistance to solvent attack and abrasion and as foams with good strength, good rebound and high impact strength.
Polyurea shows high Tg, fair resistance to greases, oils, and solvents. It can be used in truck bed liners, bridge coating, caulk and decorative designs.
Polysiloxane, siloxane-based polymers available in a wide range of physical states—from liquids to greases, waxes, resins, and rubbers. Due to perfect thermal stability (thanks to silicon, Si) uses of this material include antifoam and release agents, gaskets, seals, cable and wire insulation, hot liquids and gas conduits, etc.
Polycarbonates are transparent, self-extinguishing materials. They possess properties like crystalline thermoplasticity, high impact strength, good thermal and oxidative stability. They can be used in machinery, auto-industry, and medical applications. For example, the cockpit canopy of F-22 Raptor is made of high optical quality polycarbonate.
Polysulfides have outstanding oil and solvent resistance, good gas impermeability, good resistance to aging and ozone. However, it smells bad, and it shows low tensile strength as well as poor heat resistance. It can be used in gasoline hoses, gaskets and places that require solvent resistance and gas resistance.
Polyether shows good thermoplastic behavior, water solubility, generally good mechanical properties, moderate strength and stiffness. It is applied in sizing for cotton and synthetic fibers, stabilizers for adhesives, binders, and film formers in pharmaceuticals.
Phenol formaldehyde resin (bakelite) have good heat resistance, dimensional stability as well as good resistance to most solvents. It also shows good dielectric properties. This material is typically used in molding applications, electrical, radio, televisions and automotive parts where their good dielectric properties are of use. Some other uses include: impregnating paper, varnishes, decorative laminates for wall coverings.
Polytriazole polymers are produced from monomers which bear both an alkyne and azide functional group. The monomer units are linked to each other by the a 1,2,3-triazole group; which is produced by the 1,3-dipolar cycloaddition, also called the azide-alkyne Huisgen cycloaddition. These polymers can take on the form of a strong resin, or a gel. With oligopeptide monomers containing a terminal alkyne and terminal azide the resulting clicked peptide polymer will be biodegradable due to action of endopeptidases on the oligopeptide unit.
Branched polymers
A monomer with functionality of 3 or more will introduce branching in a polymer and will ultimately form a cross-linked macrostructure or network even at low fractional conversion. The point at which a tree-like topology transits to a network is known as the gel point because it is signalled by an abrupt change in viscosity. One of the earliest so-called thermosets is known as bakelite. It is not always water that is released in step-growth polymerization: in acyclic diene metathesis or ADMET dienes polymerize with loss of ethene.
Kinetics
The kinetics and rates of step-growth polymerization can be described using a polyesterification mechanism. The simple esterification is an acid-catalyzed process in which protonation of the acid is followed by interaction with the alcohol to produce an ester and water. However, there are a few assumptions needed with this kinetic model. The first assumption is water (or any other condensation product) is efficiently removed. Secondly, the functional group reactivities are independent of chain length. Finally, it is assumed that each step only involves one alcohol and one acid.
This is a general rate law degree of polymerization for polyesterification
where n= reaction order.
Self-catalyzed polyesterification
If no acid catalyst is added, the reaction will still proceed because the acid can act as its own catalyst. The rate of condensation at any time t can then be derived from the rate of disappearance of -COOH groups and
The second-order term arises from its use as a catalyst, and k is the rate constant. For a system with equivalent quantities of acid and glycol, the functional group concentration can be written simply as
After integration and substitution from Carothers equation, the final form is the following
For a self-catalyzed system, the number average degree of polymerization (Xn) grows proportionally with .
External catalyzed polyesterification
The uncatalyzed reaction is rather slow, and a high Xn is not readily attained. In the presence of a catalyst, there is an acceleration of the rate, and the kinetic expression is altered to
which is kinetically first order in each functional group. Hence,
and integration gives finally
For an externally catalyzed system, the number average degree of polymerization grows proportionally with .
Molecular weight distribution in linear polymerization
The product of a polymerization is a mixture of polymer molecules of different molecular weights. For theoretical and practical reasons it is of interest to discuss the distribution of molecular weights in a polymerization. The molecular weight distribution (MWD) had been derived by Flory by a statistical approach based on the concept of equal reactivity of functional groups.
Probability
Step-growth polymerization is a random process so we can use statistics to calculate the probability of finding a chain with x-structural units ("x-mer") as a function of time or conversion.
{\mathit{x}AA} + \mathit{x}BB -> AA-(BB-AA)_{\mathit{x}-1}-BB
\mathit{x}AB -> A-(B-A)_{\mathit{x}-1}-B
Probability that an 'A' functional group has reacted
Probability of finding an 'A' unreacted
Combining the above two equations leads to.
Where Px is the probability of finding a chain that is x-units long and has an unreacted 'A'. As x increases the probability decreases.
Number fraction distribution
The number fraction distribution is the fraction of x-mers in any system and equals the probability of finding it in solution.
Where N is the total number of polymer molecules present in the reaction.
Weight fraction distribution
The weight fraction distribution is the fraction of x-mers in a system and the probability of finding them in terms of mass fraction.
| Physical sciences | Organic reactions | Chemistry |
1787105 | https://en.wikipedia.org/wiki/Animal%20sexual%20behaviour | Animal sexual behaviour | Animal sexual behaviour takes many different forms, including within the same species. Common mating or reproductively motivated systems include monogamy, polygyny, polyandry, polygamy and promiscuity. Other sexual behaviour may be reproductively motivated (e.g. sex apparently due to duress or coercion and situational sexual behaviour) or non-reproductively motivated (e.g. homosexual sexual behaviour, bisexual sexual behaviour, cross-species sex, sexual arousal from objects or places, sex with dead animals, etc.).
When animal sexual behaviour is reproductively motivated, it is often termed mating or copulation; for most non-human mammals, mating and copulation occur at oestrus (the most fertile period in the mammalian female's reproductive cycle), which increases the chances of successful impregnation. Some animal sexual behaviour involves competition, sometimes fighting, between multiple males. Females often select males for mating only if they appear strong and able to protect themselves. The male that wins a fight may also have the chance to mate with a larger number of females and will therefore pass on his genes to their offspring.
Historically, it was believed that only humans and a small number of other species performed sexual acts other than for reproduction, and that animals' sexuality was instinctive and a simple "stimulus-response" behaviour. However, in addition to homosexual behaviours, a range of species masturbate and may use objects as tools to help them do so. Sexual behaviour may be tied more strongly to the establishment and maintenance of complex social bonds across a population which support its success in non-reproductive ways. Both reproductive and non-reproductive behaviours can be related to expressions of dominance over another animal or survival within a stressful situation (such as sex due to duress or coercion).
Mating systems
In sociobiology and behavioural ecology, the term "mating system" is used to describe the ways in which animal societies are structured in relation to sexual behaviour. The mating system specifies which males mate with which females, and under what circumstances. There are four basic systems:
Monogamy
Monogamy occurs when one male and one female mate exclusively with each other. A monogamous mating system is one in which individuals form long-lasting pairs and cooperate in raising offspring. These pairs may last for a lifetime, such as in pigeons, or it may occasionally change from one mating season to another, such as in emperor penguins. In contrast with tournament species, these pair-bonding species have lower levels of male aggression, competition and little sexual dimorphism. Zoologists and biologists now have evidence that monogamous pairs of animals are not always sexually exclusive. Many animals that form pairs to mate and raise offspring regularly engage in sexual activities with extra-pair partners. This includes previous examples, such as swans. Sometimes, these extra-pair sexual activities lead to offspring. Genetic tests frequently show that some of the offspring raised by a monogamous pair come from the female mating with an extra-pair male partner. These discoveries have led biologists to adopt new ways of talking about monogamy. According to Ulrich Reichard (2003):
Whatever makes a pair of animals socially monogamous does not necessarily make them sexually or genetically monogamous. Social monogamy, sexual monogamy, and genetic monogamy can occur in different combinations.
Social monogamy is relatively rare in the animal kingdom. The actual incidence of social monogamy varies greatly across different branches of the evolutionary tree. Over 90% of avian species are socially monogamous. This stands in contrast to mammals. Only 3% of mammalian species are socially monogamous, although up to 15% of primate species are. Social monogamy has also been observed in reptiles, fish, and insects.
Sexual monogamy is also rare among animals. Many socially monogamous species engage in extra-pair copulations, making them sexually non-monogamous. For example, while over 90% of birds are socially monogamous, "on average, 30% or more of the baby birds in any nest [are] sired by someone other than the resident male." Patricia Adair Gowaty has estimated that, out of 180 different species of socially monogamous songbirds, only 10% are sexually monogamous.
The incidence of genetic monogamy, determined by DNA fingerprinting, varies widely across species. For a few rare species, the incidence of genetic monogamy is 100%, with all offspring genetically related to the socially monogamous pair. But genetic monogamy is strikingly low in other species. Barash and Lipton note:
Such low levels of genetic monogamy have surprised biologists and zoologists, forcing them to rethink the role of social monogamy in evolution. They can no longer assume social monogamy determines how genes are distributed in a species. The lower the rates of genetic monogamy among socially monogamous pairs, the less of a role social monogamy plays in determining how genes are distributed among offspring.
Polygamy
The term polygamy is an umbrella term used to refer generally to non-monogamous matings. As such, polygamous relationships can be polygynous, polyandrous or polygynandrous. In a small number of species, individuals can display either polygamous or monogamous behaviour depending on environmental conditions. An example is the social wasp Apoica flavissima. In some species, polygyny and polyandry is displayed by both sexes in the population. Polygamy in both sexes has been observed in red flour beetle (Tribolium castaneum). Polygamy is also seen in many Lepidoptera species including Mythimna unipuncta (true armyworm moth).
A tournament species is one in which "mating tends to be highly polygamous and involves high levels of male-male aggression and competition." Tournament behaviour often correlates with high levels of sexual dimorphism, examples of species including chimpanzees and baboons. Most polygamous species present high levels of tournament behaviour, with a notable exception being bonobos.
Polygyny
Polygyny occurs when one male gets exclusive mating rights with multiple females. In some species, notably those with harem-like structures, only one of a few males in a group of females will mate. Technically, polygyny in sociobiology and zoology is defined as a system in which a male has a relationship with more than one female, but the females are predominantly bonded to a single male. Should the active male be driven out, killed, or otherwise removed from the group, in a number of species the new male will ensure that breeding resources are not wasted on another male's young. The new male may achieve this in many different ways, including:
competitive infanticide: in lions, hippopotamuses, and some monkeys, the new male will kill the offspring of the previous alpha male to cause their mothers to become receptive to his sexual advances since they are no longer nursing. To prevent this, many female primates exhibit ovulation cues among all males, and show situation-dependent receptivity.
harassment to miscarriage: amongst wild horses and baboons, the male will continually attack pregnant females until they miscarry.
Pheromone-based spontaneous abortion
in some rodents such as mice, a new male with a different scent will cause females who are pregnant to spontaneously fail to implant recently fertilised eggs. This does not require contact; it is mediated by scent alone. It is known as the Bruce effect.
Von Haartman specifically described the mating behaviour of the European pied flycatcher as successive polygyny. Within this system, the males leave their home territory once their primary female lays her first egg. Males then create a second territory, presumably in order to attract a secondary female to breed. Even when they succeed at acquiring a second mate, the males typically return to the first female to exclusively provide for her and her offspring.
Polygynous mating structures are estimated to occur in up to 90% of mammal species. As polygyny is the most common form of polygamy among vertebrates (including humans), it has been studied far more extensively than polyandry or polygynandry.
Polyandry
Polyandry occurs when one female gets exclusive mating rights with multiple males. In some species, such as redlip blennies, both polygyny and polyandry are observed.
The males in some deep sea anglerfishes are much smaller than the females. When they find a female they bite into her skin, releasing an enzyme that digests the skin of their mouths and her body and fusing the pair down to the blood-vessel level. The male then slowly atrophies, losing first his digestive organs, then his brain, heart, and eyes, ending as nothing more than a pair of gonads, which release sperm in response to hormones in the female's bloodstream indicating egg release. This extreme sexual dimorphism ensures that, when the female is ready to spawn, she has a mate immediately available. A single anglerfish female can "mate" with many males in this manner.
Polygynandry
Polygynandry occurs when multiple males mate indiscriminately with multiple females. The numbers of males and females need not be equal, and in vertebrate species studied so far, there are usually fewer males. Two examples of systems in primates are promiscuous mating chimpanzees and bonobos. These species live in social groups consisting of several males and several females. Each female copulates with many males, and vice versa. In bonobos, the amount of promiscuity is particularly striking because bonobos use sex to alleviate social conflict as well as to reproduce. This mutual promiscuity is the approach most commonly used by spawning animals, and is perhaps the "original fish mating system." Common examples are forage fish, such as herrings, which form huge mating shoals in shallow water. The water becomes milky with sperm and the bottom is draped with millions of fertilised eggs.
Parental investment and reproductive success
Female and male sexual behaviour differ in many species. Often, males are more active in initiating mating, and bear the more conspicuous sexual ornamentation like antlers and colourful plumage. This is a result of anisogamy, where sperm are smaller and much less costly (energetically) to produce than eggs. This difference in physiological cost means that males are more limited by the number of mates they can secure, while females are limited by the quality of genes of her mates, a phenomenon known as Bateman's principle. Many females also have extra reproductive burdens in that parental care often falls mainly, or exclusively, on them. Thus, females are more limited in their potential reproductive success. In species where males take on more of the reproductive costs, such as sea horses and jacanas, the role is reversed, and the females are larger, more aggressive and more brightly coloured than the males.
In hermaphroditic animals, the costs of parental care can be evenly distributed between the sexes, e.g. earthworms. In some species of planarians, sexual behaviour takes the form of penis fencing. In this form of copulation, the individual that first penetrates the other with the penis, forces the other to be female, thus carrying the majority of the cost of reproduction. Post mating, banana slugs will some times gnaw off their partners penis as an act of sperm competition called apophallation. This is costly as they must heal, and spend more energy courting conspecifics that can act as male and female. A hypothesis suggests these slugs may be able to compensate the loss of the male function by directing energy that would have been put towards it to the female function. In the grey slug, the sharing of cost leads to a spectacular display, where the mates suspend themselves high above the ground from a slime thread, ensuring none of them can refrain from taking on the cost of egg-bearer.
Seasonality
Many animal species have specific mating (or breeding) periods e.g. (seasonal breeding) so that offspring are born or hatch at an optimal time. In marine species with limited mobility and external fertilisation like corals, sea urchins and clams, the timing of the common spawning is the only externally visible form of sexual behaviour. In areas with continuously high primary production, some species have a series of breeding seasons throughout the year. This is the case with most primates (who are primarily tropical and subtropical animals). Some animals (opportunistic breeders) breed dependent upon other conditions in their environment aside from time of year.
Mammals
Mating seasons are often associated with changes to herd or group structure, and behavioural changes, including territorialism amongst individuals. These may be annual (e.g. wolves), biannual (e.g. dogs) or more frequently (e.g. horses). During these periods, females of most mammalian species are more mentally and physically receptive to sexual advances, a period scientifically described as oestrus but commonly described as being "in season" or "in heat". Sexual behaviour may occur outside oestrus, and such acts as do occur are not necessarily harmful.
Some mammals (e.g. domestic cats, rabbits and camelids) are termed "induced ovulators". For these species, the female ovulates due to an external stimulus during, or just prior to, mating, rather than ovulating cyclically or spontaneously. Stimuli causing induced ovulation include the sexual behaviour of coitus, sperm and pheromones. Domestic cats have penile spines. Upon withdrawal of a cat's penis, the spines rake the walls of the female's vagina, which may cause ovulation.
Amphibians
For many amphibians, an annual breeding cycle applies, typically regulated by ambient temperature, precipitation, availability of surface water and food supply. This breeding season is accentuated in temperate regions, in boreal climate the breeding season is typically concentrated to a few short days in the spring. Some species, such as the Rana clamitans (green frog), spend from June to August defending their territory. In order to protect these territories, they use five vocalizations.
Fish
Like many coral reef dwellers, the clownfish spawn around the time of the full moon in the wild. In a group of clownfish, there is a strict dominance hierarchy. The largest and most aggressive female is found at the top. Only two clownfish, a male and a female, in a group reproduce through external fertilisation. Clownfish are sequential hermaphrodites, meaning that they develop into males first, and when they mature, they become females. If the female clownfish is removed from the group, such as by death, one of the largest and most dominant males will become a female. The remaining males will move up a rank in the hierarchy.
Motivation
Various neurohormones stimulate sexual wanting in animals. In general, studies have suggested that dopamine is involved in sexual incentive motivation, oxytocin and melanocortins in sexual attraction, and noradrenaline in sexual arousal. Vasopressin is also involved in the sexual behaviour of some animals.
Neurohormones in the mating systems of voles
The mating system of prairie voles is monogamous; after mating, they form a lifelong bond. In contrast, montane voles have a polygamous mating system. When montane voles mate, they form no strong attachments, and separate after copulation. Studies on the brains of these two species have found that it is two neurohormones and their respective receptors that are responsible for these differences in mating strategies. Male prairie voles release vasopressin after copulation with a partner, and an attachment to their partner then develops. Female prairie voles release oxytocin after copulation with a partner, and similarly develop an attachment to their partner.
Neither male nor female montane voles release high quantities of oxytocin or vasopressin when they mate. Even when injected with these neurohormones, their mating system does not change. In contrast, if prairie voles are injected with the neurohormones, they may form a lifelong attachment, even if they have not mated. The differing response to the neurohormones between the two species is due to a difference in the number of oxytocin and vasopressin receptors. Prairie voles have a greater number of oxytocin and vasopressin receptors compared to montane voles, and are therefore more sensitive to those two neurohormones. It's believed that it's the quantity of receptors, rather than the quantity of the hormones, that determines the mating system and bond-formation of either species.
Oxytocin and rat sexual behaviour
Mother rats experience a postpartum estrus which makes them highly motivated to mate. However, they also have a strong motivation to protect their newly born pups. As a consequence, the mother rat solicits males to the nest but simultaneously becomes aggressive towards them to protect her young. If the mother rat is given injections of an oxytocin receptor antagonist, they no longer experience these maternal motivations.
Prolactin influences social bonding in rats.
Oxytocin and primate sexual behaviour
Oxytocin plays a similar role in non-human primates as it does in humans.
Grooming, sex, and cuddling frequencies correlate positively with levels of oxytocin. As the level of oxytocin increases so does sexual motivation. While oxytocin plays a major role in parent child relationships, it is also found to play a role in adult sexual relationships. Its secretion affects the nature of the relationship or if there will even be a relationship at all.
Studies have shown that oxytocin is higher in monkeys in lifelong monogamous relationships compared to monkeys which are single. Furthermore, the oxytocin levels of the couples correlate positively; when the oxytocin secretion of one increases, the other one also increases. Higher levels of oxytocin are related to monkeys expressing more behaviours such as cuddling, grooming and sex, while lower levels of oxytocin reduce motivation for these activities.
Research on oxytocin's role in the animal brain suggests that it plays less of a role in behaviours of love and affection than previously believed. "When oxytocin was first discovered in 1909, it was thought mostly to influence a mother's labour contractions and milk let-down. Then, in the 1990s, research with prairie voles found that giving them a dose of oxytocin resulted in the formation of a bond with their future mate (Azar, 40)." Oxytocin has since been treated by the media as the sole player in the "love and mating game" in mammals. This view, however, is proving to be false as, "most hormones don't influence behaviour directly. Rather, they affect thinking and emotions in variable ways (Azar, 40)." There is much more involved in sexual behaviour in the mammalian animal than oxytocin and vasopressin can explain.
Pleasure
It is often assumed that animals do not have sex for pleasure, or alternatively that humans, pigs, bonobos (and perhaps dolphins and one or two more species of primates) are the only species that do. This is sometimes stated as "animals mate only for reproduction". This view is considered a misconception by some scholars. Jonathan Balcombe argues that the prevalence of non-reproductive sexual behaviour in certain species suggests that sexual stimulation is pleasurable. He also points to the presence of the clitoris in some female mammals, and evidence for female orgasm in primates. On the other hand, it is impossible to know the subjective feelings of animals, and the notion that non-human animals experience emotions similar to humans is a contentious subject.
A 2006 Danish Animal Ethics Council report, which examined current knowledge of animal sexuality in the context of legal queries concerning sexual acts by humans, has the following comments, primarily related to domestically common animals:
Koinophilia
Koinophilia is the love of the "normal" or phenotypically common (from the Greek , , meaning "usual" or "common"). The term was introduced to scientific literature in 1990, and refers to the tendency of animals seeking a mate to prefer that mate not to have any unusual, peculiar or deviant features. Similarly, animals preferentially choose mates with low fluctuating asymmetry. However, animal sexual ornaments can evolve through runaway selection, which is driven by (usually female) selection for non-standard traits.
Interpretation bias
The field of study of sexuality in non-human species was a long-standing taboo. In the past, researchers sometimes failed to observe, miscategorised or misdescribed sexual behaviour which did not meet their preconceptions—their bias tended to support what would now be described as conservative sexual mores. An example of overlooking behaviour relates to descriptions of giraffe mating:
In the 21st century, liberal social or sexual views are often projected upon animal subjects of research. Popular discussions of bonobos are a frequently cited example. Current research frequently expresses views such as that of the Natural History Museum at the University of Oslo, which in 2006 held an exhibition on animal sexuality:
Other animal activities may be misinterpreted due to the frequency and context in which animals perform the behaviour. For example, domestic ruminants display behaviours such as mounting and head-butting. This often occurs when the animals are establishing dominance relationships and are not necessarily sexually motivated. Careful analysis must be made to interpret what animal motivations are being expressed by those behaviours.
Types of sexual behaviour
Reproductive sexual behaviour
Copulation
Copulation is the union of the male and female sex organs, the sexual activity specifically organized to transmit male sperm into the body of the female.
Cuckoldry
Alternative male strategies which allow small males to engage in cuckoldry can develop in species such as fish where spawning is dominated by large and aggressive males. Cuckoldry is a variant of polyandry, and can occur with sneak spawners. A sneak spawner is a male that rushes in to join the spawning rush of a spawning pair. A spawning rush occurs when a fish makes a burst of speed, usually on a near vertical incline, releasing gametes at the apex, followed by a rapid return to the lake or sea floor or fish aggregation. Sneaking males do not take part in courtship. In salmon and trout, for example, jack males are common. These are small silvery males that migrate upstream along with the standard, large, hook-nosed males and that spawn by sneaking into redds to release sperm simultaneously with a mated pair. This behaviour is an evolutionarily stable strategy for reproduction, because it is favoured by natural selection just like the "standard" strategy of large males.
Hermaphroditism
Hermaphroditism occurs when a given individual in a species possesses both male and female reproductive organs, or can alternate between possessing first one, and then the other. Hermaphroditism is common in invertebrates but rare in vertebrates. It can be contrasted with gonochorism, where each individual in a species is either male or female, and remains that way throughout their lives. Most fish are gonochorists, but hermaphroditism is known to occur in 14 families of teleost fishes.
Usually hermaphrodites are sequential, meaning they can switch sex, usually from female to male (protogyny). This can happen if a dominant male is removed from a group of females. The largest female in the harem can switch sex over a few days and replace the dominant male. This is found amongst coral reef fishes such as groupers, parrotfishes and wrasses. As an example, most wrasses are protogynous hermaphrodites within a haremic mating system. It is less common for a male to switch to a female (protandry). A common example of a protandrous species are clownfish—if the larger, dominant female dies, in many cases, the reproductive male gains weight and becomes the female. Hermaphroditism allows for complex mating systems. Wrasses exhibit three different mating systems: polygynous, lek-like, and promiscuous mating systems.
Sexual cannibalism
Sexual cannibalism is a behaviour in which a female animal kills and consumes the male before, during, or after copulation. Sexual cannibalism confers fitness advantages to both the male and female. Sexual cannibalism is common among insects, arachnids and amphipods. There is also evidence of sexual cannibalism in gastropods and copepods.
Sexual coercion
Sex in a forceful or apparently coercive context has been documented in a variety of species. In some herbivorous herd species, or species where males and females are very different in size, the male dominates sexually by force and size.
Some species of birds have been observed combining sexual intercourse with apparent violent assault; these include ducks, and geese. Female white-fronted bee-eaters are subjected to forced copulations. When females emerge from their nest burrows, males sometimes force them to the ground and mate with them. Such forced copulations are made preferentially on females who are laying and who may therefore lay eggs fertilized by the male.
It has been reported that young male elephants in South Africa sexually coerced and killed rhinoceroses. This interpretation of the elephants' behaviour was disputed by one of the original study's authors, who said there was "nothing sexual about these attacks".
Parthenogenesis
Parthenogenesis is a form of asexual reproduction in which growth and development of embryos occur without fertilisation. Technically, parthenogenesis is not a behaviour, however, sexual behaviours may be involved.
Whip-tailed lizard females have the ability to reproduce through parthenogenesis and as such males are rare and sexual breeding non-standard. Females engage in "pseudocopulation" to stimulate ovulation, with their behaviour following their hormonal cycles; during low levels of oestrogen, these (female) lizards engage in "masculine" sexual roles. Those animals with currently high oestrogen levels assume "feminine" sexual roles. Lizards that perform the courtship ritual have greater fecundity than those kept in isolation due to an increase in hormones triggered by the sexual behaviours. So, even though asexual whiptail lizards populations lack males, sexual stimuli still increase reproductive success. From an evolutionary standpoint these females are passing their full genetic code to all of their offspring rather than the 50% of genes that would be passed in sexual reproduction.
It is rare to find true parthenogenesis in fishes, where females produce female offspring with no input from males. All-female species include the Texas silverside, Menidia clarkhubbsi and a complex of Mexican mollies.
Parthenogenesis has been recorded in 70 vertebrate species including hammerhead sharks, blacktip sharks, amphibians
and lizards.
Unisexuality
Unisexuality occurs when a species is all-male or all-female. Unisexuality occurs in some fish species and can take complex forms. Squalius alburnoides, a minnow found in several river basins in Portugal and Spain, appears to be an all-male species. The existence of this species illustrates the potential complexity of mating systems in fish. The species originated as a hybrid between two species and is diploid but not hermaphroditic. It can have triploid and tetraploid forms, including all-female forms that reproduce mainly through hybridogenesis.
Others
Interbreeding: Hybrid offspring can result from the mating of two organisms of distinct but closely related parent species, although the resulting offspring is not always fertile. According to Alfred Kinsey, genetic studies on wild animal populations have shown a "large number" of inter-species hybrids.
Prostitution: There are reports that animals occasionally engage in prostitution. A small number of pair-bonded females within a group of penguins took nesting material (stones) after copulating with a non-partner male. The researcher stated "I was watching opportunistically, so I can't give an exact figure of how common it really is." It has been reported that "bartering of meat for sex ... forms part of the social fabric of a troop of wild chimps living in the Tai National Park in the Côte d'Ivoire."
Pavlovian conditioning: The sexualisation of objects or locations is recognised in the animal breeding world. For example, male animals may become sexually aroused upon visiting a location where they have been allowed to have sex before, or upon seeing a stimulus previously associated with sexual activity such as an artificial vagina. Sexual preferences for certain cues can be artificially induced in rats by pairing scents or objects with their early sexual experiences. The primary motivation of this behaviour is Pavlovian conditioning, and the association is due to a conditioned response (or association) formed with a distinctive "reward".
Viewing images: A study using four adult male rhesus macaques (Macaca mulatta) showed that male rhesus macaques will give up a highly valued item, juice, to see images of the faces or perineum of high-status females. Encouraging captive pandas to mate is problematic. Showing young male pandas "panda pornography" is credited with a recent population boom among pandas in captivity in China. One researcher attributed the success to the sounds on the recordings.
Copulatory wounding and traumatic insemination: Injury to a partner's genital tract during mating occurs in at least 40 taxa, ranging from fruit flies to humans. However, it often goes unnoticed due to its cryptic nature and because of internal wounds not visible outside.
Non-reproductive sexual behaviour
There is a range of behaviours that animals perform that appear to be sexually motivated but which can not result in reproduction. These include:
Masturbation: Some species, both male and female, masturbate, both when partners are available and otherwise.
Oral sex: Several species engage in both autofellatio and oral sex. This has been documented in brown bears, Tibetan macaques, wolves, goats, primates, bats, cape ground squirrels and sheep. In the greater short-nosed fruit bat, copulation by males is dorsoventral and the females lick the shaft or the base of the male's penis, but not the glans which has already penetrated the vagina. While the females do this, the penis is not withdrawn and research has shown a positive relationship between length of the time that the penis is licked and the duration of copulation. Post copulation genital grooming has also been observed.
Homosexuality: Same-sex sexual behaviour occurs in a range of species, especially in social species, particularly in marine birds and mammals, monkeys, and the great apes. , the scientific literature contained reports of homosexual behaviour in at least 471 wild species. Organisers of the Against Nature? exhibit stated that "homosexuality has been observed among 1,500 species, and that in 500 of those it is well documented."
Genital-genital rubbing: This is sexual activity in which one animal rubs his or her genitals against the genitals of another animal. This is stated to be the "bonobo's most typical sexual pattern, undocumented in any other primate".
Inter-species mating: Some animals opportunistically mate with individuals of another species.
Sex involving juveniles: Male stoats (Mustela erminea) will sometimes mate with infant females of their species. This is a natural part of their reproductive biology—they have a delayed gestation period, so these females give birth the following year when they are fully grown. Juvenile male common chimpanzees have been recorded mounting and copulating with immature chimps. Infants in bonobo societies are often involved in sexual behaviour.
Necrophilia: This describes when an animal engages in a sexual act with a dead animal. It has been observed in mammals, birds, reptiles and frogs.
Bisexuality: This describes when an animal shows sexual behaviour towards both males and females.
Extended female sexuality: This is when females mate with males outside of their conceptive period.
Seahorse
Seahorses, once considered to be monogamous species with pairs mating for life, were described in a 2007 study as "promiscuous, flighty, and more than a little bit gay". Scientists at 15 aquaria studied 90 seahorses of three species. Of 3,168 sexual encounters, 37% were same-sex acts. Flirting was common (up to 25 potential partners a day of both sexes); only one species (the British spiny seahorse) included faithful representatives, and for these 5 of 17 were faithful, 12 were not. Bisexual behaviour was widespread and considered "both a great surprise and a shock", with big-bellied seahorses of both sexes not showing partner preference. 1,986 contacts were male-female, 836 were female-female and 346 were male-male.
Bonobo
Among bonobos, males and females engage in sexual behaviour with the same and the opposite sex, with females being particularly noted for engaging in sexual behaviour with each other and at up to 75% of sexual activity being non-reproductive, as being sexually active does not necessarily correlate with their ovulation cycles. Sexual activity occurs between almost all ages and sexes of bonobo societies. Primatologist Frans de Waal believes that bonobos use sexual activity to resolve conflict between individuals. Immature bonobos, contrariwise, perform genital contact when relaxed.
Macaque
Similar same-sex sexual behaviours occur in both male and female macaques. It is thought to be done for pleasure as an erect male mounts and thrusts upon or into another male. Sexual receptivity can also be indicated by red faces and shrieking. Mutual ejaculation after a combination of anal intercourse and masturbation has also been witnessed, although it may be rare. In comparison to socio-sexual behaviours such as dominance displays, homosexual mounts last longer, happen in series, and usually involve pelvic thrusting.
Females are also thought to participate for pleasure as vulvar, perineal, and anal stimulation is part of these interactions. The stimulation can come from their own tails, mounting their partner, thrusting or a combination of these.
Dolphin
Male bottlenose dolphins have been observed working in pairs to follow or restrict the movement of a female for weeks at a time, waiting for her to become sexually receptive. The same pairs have also been observed engaging in intense sexual play with each other. Janet Mann, a professor of biology and psychology at Georgetown University, argues that the common same-sex behaviour among male dolphin calves is about bond formation and benefits the species evolutionarily. Studies have shown the dolphins later in life are bisexual and the male bonds forged from homosexuality work for protection as well as locating females with which to reproduce.
In 1991, an English man was prosecuted for allegedly having sexual contact with a dolphin. The man was found not guilty after it was revealed at trial that the dolphin was known to tow bathers through the water by hooking his penis around them.
Hyena
The female spotted hyena has a unique urinary-genital system, closely resembling the penis of the male, called a pseudo-penis. Dominance relationships with strong sexual elements are routinely observed between related females. They are notable for using visible sexual arousal as a sign of submission but not dominance in males as well as females (females have a sizeable erectile clitoris). It is speculated that to facilitate this, their sympathetic and parasympathetic nervous systems may be partially reversed in respect to their reproductive organs.
Mating behaviour
Vertebrates
Mammals
Mammals mate by vaginal copulation. To achieve this, the male usually mounts the female from behind. The female may exhibit lordosis in which she arches her back ventrally to facilitate entry of the penis, which is particularly present in elephants, felids, and rodents. Amongst the land mammals, other than humans, only bonobos mate in a face-to-face position, as the females' anatomy seems to reflect, although ventro-ventral copulation has also been observed in Rhabdomys. Some sea mammals copulate in a belly-to-belly position. Some camelids mate in a lying-down position. In most mammals ejaculation occurs after multiple intromissions, but in most primates, copulation consists of one brief intromission. In most ruminant species, a single pelvic thrust occurs during copulation. In most deer species, a copulatory jump also occurs.
During mating, a "copulatory tie" occurs in mammals such as fossas, canids with the exception of African wild dogs, and Japanese martens. A "copulatory lock" also occurs in some primate species, such as Galago senegalensis.
The copulatory behaviour of many mammalian species is affected by sperm competition.
Some females have concealed fertility, making it difficult for males to evaluate if a female is fertile - humans are amongst these species. This is costly as ejaculation expends much energy.
Invertebrates
Invertebrates are often hermaphrodites. Some hermaphroditic land snails begin mating with an elaborate tactile courting ritual. The two snails circle around each other for up to six hours, touching with their tentacles, and biting lips and the area of the genital pore, which shows some preliminary signs of the eversion of the penis. As the snails approach mating, hydraulic pressure builds up in the blood sinus surrounding an organ housing a sharpened dart. The dart is made of calcium carbonate or chitin, and is called a love dart. Each snail manoeuvres to get its genital pore in the best position, close to the other snail's body. Then, when the body of one snail touches the other snail's genital pore, it triggers the firing of the love dart. After the snails have fired their darts, they copulate and exchange sperm as a separate part of the mating progression. The love darts are covered with a mucus that contains a hormone-like substance that facilitates the survival of the sperm.
Penis fencing is a mating behaviour engaged in by certain species of flatworm, such as Pseudobiceros bedfordi. Species which engage in the practice are hermaphroditic, possessing both eggs and sperm-producing testes. The species "fence" using two-headed dagger-like penises which are pointed, and white in colour. One organism inseminates the other. The sperm is absorbed through pores in the skin, causing fertilisation.
Corals can be both gonochoristic (unisexual) and hermaphroditic, each of which can reproduce sexually and asexually. Reproduction also allows corals to settle new areas. Corals predominantly reproduce sexually. 25% of hermatypic corals (stony corals) form single sex (gonochoristic) colonies, while the rest are hermaphroditic. About 75% of all hermatypic corals "broadcast spawn" by releasing gametes eggs and sperm into the water to spread offspring. The gametes fuse during fertilisation to form a microscopic larva called a planula, typically pink and elliptical in shape. Synchronous spawning is very typical on the coral reef and often, even when multiple species are present, all corals spawn on the same night. This synchrony is essential so that male and female gametes can meet. Corals must rely on environmental cues, varying from species to species, to determine the proper time to release gametes into the water. The cues involve lunar changes, sunset time, and possibly chemical signalling. Synchronous spawning may form hybrids and is perhaps involved in coral speciation.
Butterflies spend much time searching for mates. When the male spots a mate, he will fly closer and release pheromones. He then performs a special courtship dance to attract the female. If the female appreciates the dancing she may join him. Then they join their bodies together end to end at their abdomens. Here, the male passes the sperm to the female's egg-laying tube, which will soon be fertilised by the sperm.
Many animals make plugs of mucus to seal the female's orifice after mating. Normally such plugs are secreted by the male, to block subsequent partners. In spiders the female can assist the process. Spider sex is unusual in that males transfer their sperm to the female on small limbs called pedipalps. They use these to pick their sperm up from their genitals and insert it into the female's sexual orifice, rather than copulating directly. On the 14 occasions a sexual plug was made, the female produced it without assistance from the male. On ten of these occasions the male's pedipalps then seemed to get stuck while he was transferring the sperm (which is rarely the case in other species of spider), and he had great difficulty freeing himself. In two of those ten instances, he was eaten as a result.
In the orb-weaving spider species Zygiella x-notata, individuals engage in a variety of sexual behaviors including male choosiness, mate guarding, and vibrational signaling in courtship.
Genetic evidence of interspecies sexual activity in humans
Research into human evolution confirms that, in some cases, interspecies sexual activity may have been responsible for the evolution of new species (speciation). Analysis of animal genes found evidence that, after humans had diverged from other apes, interspecies mating nonetheless occurred regularly enough to change certain genes in the new gene pool. Researchers found that the X chromosomes of humans and chimps may have diverged around 1.2 million years after the other chromosomes. One possible explanation is that modern humans emerged from a hybrid of human and chimp populations. A 2012 study questioned this explanation, concluding that "there is no strong reason to involve complicated factors in explaining the autosomal data".
Inbreeding avoidance
When close relatives mate, progeny may exhibit the detrimental effects of inbreeding depression. Inbreeding depression is predominantly caused by the homozygous expression of recessive deleterious alleles. Over time, inbreeding depression may lead to the evolution of inbreeding avoidance behaviour. Several examples of animal behaviour that reduce mating of close relatives and inbreeding depression are described next.
Reproductively active female naked mole-rats tend to associate with unfamiliar males (usually non-kin), whereas reproductively inactive females do not discriminate. The preference of reproductively active females for unfamiliar males is interpreted as an adaptation for avoiding inbreeding.
When mice inbreed with close relatives in their natural habitat, there is a significant detrimental effect on progeny survival. In the house mouse, the major urinary protein (MUP) gene cluster provides a highly polymorphic scent signal of genetic identity that appears to underlie kin recognition and inbreeding avoidance. Thus there are fewer matings between mice sharing MUP haplotypes than would be expected if there were random mating.
Meerkat females appear to be able to discriminate the odour of their kin from the odour of their non-kin. Kin recognition is a useful ability that facilitates both cooperation among relatives and the avoidance of inbreeding. When mating does occur between meerkat relatives, it often results in inbreeding depression. Inbreeding depression was evident for a variety of traits: pup mass at emergence from the natal burrow, hind-foot length, growth until independence and juvenile survival.
The grey-sided vole (Myodes rufocanus) exhibits male-biased dispersal as a means of avoiding incestuous matings. Among those matings that do involve inbreeding the number of weaned juveniles in litters is significantly smaller than that from non-inbred litters indicating inbreeding depression.
In natural populations of the bird Parus major (great tit), inbreeding is likely avoided by dispersal of individuals from their birthplace, which reduces the chance of mating with a close relative. Dispersing to avoid inbreeding is a common behavior amongst animals, such as felids and canids, although inbreeding can still occur, albeit rarely.
Toads display breeding site fidelity, as do many amphibians. Individuals that return to natal ponds to breed will likely encounter siblings as potential mates. Although incest is possible, Bufo americanus siblings rarely mate. These toads likely recognise and actively avoid close kins as mates. Advertisement vocalisations by males appear to serve as cues by which females recognise their kin.
| Biology and health sciences | Ethology | Biology |
1787246 | https://en.wikipedia.org/wiki/Ion%20chromatography | Ion chromatography | Ion chromatography (or ion-exchange chromatography) is a form of chromatography that separates ions and ionizable polar molecules based on their affinity to the ion exchanger. It works on almost any kind of charged molecule—including small inorganic anions, large proteins, small nucleotides, and amino acids. However, ion chromatography must be done in conditions that are one pH unit away from the isoelectric point of a protein.
The two types of ion chromatography are anion-exchange and cation-exchange. Cation-exchange chromatography is used when the molecule of interest is positively charged. The molecule is positively charged because the pH for chromatography is less than the pI (also known as pH(I)). In this type of chromatography, the stationary phase is negatively charged and positively charged molecules are loaded to be attracted to it. Anion-exchange chromatography is when the stationary phase is positively charged and negatively charged molecules (meaning that pH for chromatography is greater than the pI) are loaded to be attracted to it. It is often used in protein purification, water analysis, and quality control. The water-soluble and charged molecules such as proteins, amino acids, and peptides bind to moieties which are oppositely charged by forming ionic bonds to the insoluble stationary phase. The equilibrated stationary phase consists of an ionizable functional group where the targeted molecules of a mixture to be separated and quantified can bind while passing through the column—a cationic stationary phase is used to separate anions and an anionic stationary phase is used to separate cations. Cation exchange chromatography is used when the desired molecules to separate are cations and anion exchange chromatography is used to separate anions. The bound molecules then can be eluted and collected using an eluant which contains anions and cations by running a higher concentration of ions through the column or by changing the pH of the column.
One of the primary advantages for the use of ion chromatography is that only one interaction is involved the separation, as opposed to other separation techniques; therefore, ion chromatography may have higher matrix tolerance. Another advantage of ion exchange is the predictability of elution patterns (based on the presence of the ionizable group). For example, when cation exchange chromatography is used, certain cations will elute out first and others later. A local charge balance is always maintained. However, there are also disadvantages involved when performing ion-exchange chromatography, such as constant evolution of the technique which leads to the inconsistency from column to column. A major limitation to this purification technique is that it is limited to ionizable group.
History
Ion chromatography has advanced through the accumulation of knowledge over a course of many years. Starting from 1947, Spedding and Powell used displacement ion-exchange chromatography for the separation of the rare earths. Additionally, they showed the ion-exchange separation of 14N and 15N isotopes in ammonia. At the start of the 1950s, Kraus and Nelson demonstrated the use of many analytical methods for metal ions dependent on their separation of their chloride, fluoride, nitrate or sulfate complexes by anion chromatography. Automatic in-line detection was progressively introduced from 1960 to 1980 as well as novel chromatographic methods for metal ion separations. A groundbreaking method by Small, Stevens and Bauman at Dow Chemical Co. unfolded the creation of the modern ion chromatography. Anions and cations could now be separated efficiently by a system of suppressed conductivity detection. In 1979, a method for anion chromatography with non-suppressed conductivity detection was introduced by Gjerde et al. Following it in 1980, was a similar method for cation chromatography.
As a result, a period of extreme competition began within the IC market, with supporters for both suppressed and non-suppressed conductivity detection. This competition led to fast growth of new forms and the fast evolution of IC. A challenge that needs to be overcome in the future development of IC is the preparation of highly efficient monolithic ion-exchange columns and overcoming this challenge would be of great importance to the development of IC.
The boom of Ion exchange chromatography primarily began between 1935 and 1950 during World War II and it was through the "Manhattan project" that applications and IC were significantly extended. Ion chromatography was originally introduced by two English researchers, agricultural Sir Thompson and chemist J T Way. The works of Thompson and Way involved the action of water-soluble fertilizer salts, ammonium sulfate and potassium chloride. These salts could not easily be extracted from the ground due to the rain. They performed ion methods to treat clays with the salts, resulting in the extraction of ammonia in addition to the release of calcium. It was in the fifties and sixties that theoretical models were developed for IC for further understanding and it was not until the seventies that continuous detectors were utilized, paving the path for the development from low-pressure to high-performance chromatography. Not until 1975 was "ion chromatography" established as a name in reference to the techniques, and was thereafter used as a name for marketing purposes. Today IC is important for investigating aqueous systems, such as drinking water. It is a popular method for analyzing anionic elements or complexes that help solve environmentally relevant problems. Likewise, it also has great uses in the semiconductor industry.
Because of the abundant separating columns, elution systems, and detectors available, chromatography has developed into the main method for ion analysis.
When this technique was initially developed, it was primarily used for water treatment. Since 1935, ion exchange chromatography rapidly manifested into one of the most heavily leveraged techniques, with its principles often being applied to majority of fields of chemistry, including distillation, adsorption, and filtration.
Principle
Ion-exchange chromatography separates molecules based on their respective charged groups. Ion-exchange chromatography retains analyte molecules on the column based on coulombic (ionic) interactions. The ion exchange chromatography matrix consists of positively and negatively charged ions. Essentially, molecules undergo electrostatic interactions with opposite charges on the stationary phase matrix. The stationary phase consists of an immobile matrix that contains charged ionizable functional groups or ligands. The stationary phase surface displays ionic functional groups (R-X) that interact with analyte ions of opposite charge. To achieve electroneutrality, these immobilized charges couple with exchangeable counterions in the solution. Ionizable molecules that are to be purified, compete with these exchangeable counterions, for binding to the immobilized charges on the stationary phase. These ionizable molecules are retained or eluted based on their charge. Initially, molecules that do not bind or bind weakly to the stationary phase are first to be washed away. Altered conditions are needed for the elution of the molecules that bind to the stationary phase. The concentration of the exchangeable counterions, which competes with the molecules for binding, can be increased, or the pH can be changed to affect the ionic charge of the eluent or the solute. A change in pH affects the charge on the particular molecules and, therefore, alter their binding. When reducing the net charge of the solute's molecules, they start eluting out. This way, such adjustments can be used to release the proteins of interest. Additionally, concentration of counterions can be gradually varied to affect the retention of the ionized molecules, thus separate them. This type of elution is called gradient elution. On the other hand, step elution can be used, in which the concentration of counterions are varied in steps. This type of chromatography is further subdivided into cation exchange chromatography and anion-exchange chromatography. Positively charged molecules bind to cation exchange resins, while negatively charged molecules bind to anion exchange resins. The ionic compound consisting of the cationic species M+ and the anionic species B− can be retained by the stationary phase.
Cation exchange chromatography retains positively charged cations because the stationary phase displays a negatively charged functional group:
Anion exchange chromatography retains anions using positively charged functional group:
Note that the ion strength of either C+ or A− in the mobile phase can be adjusted to shift the equilibrium position, thus retention time.
The ion chromatogram shows a typical chromatogram obtained with an anion exchange column.
Procedure
Before ion-exchange chromatography can be initiated, it must be equilibrated. The stationary phase must be equilibrated to certain requirements that depend on the experiment that you are working with. Once equilibrated, the charged ions in the stationary phase will be attached to its opposite charged exchangeable ions, such as Cl− or Na+. Next, a buffer should be chosen in which the desired protein can bind to. After equilibration, the column needs to be washed. The washing phase will help elute out all impurities that does not bind to the matrix while the protein of interest remains bounded. This sample buffer needs to have the same pH as the buffer used for equilibration to help bind the desired proteins. Uncharged proteins will be eluted out of the column at a similar speed of the buffer flowing through the column with no retention. Once the sample has been loaded onto to the column, and the column has been washed with the buffer to elute out all non-desired proteins, elution is carried out at specific conditions to elute the desired proteins that are bound to the matrix. Bound proteins are eluted out by utilizing a gradient of linearly increasing salt concentration. With increasing ionic strength of the buffer, the salt ions will compete with the desired proteins in order to bind to charged groups on the surface of the medium. This will cause desired proteins to be eluted out of the column. Proteins that have a low net charge will be eluted out first as the salt concentration increases causing the ionic strength to increase. Proteins with high net charge will need a higher ionic strength for them to be eluted out of the column.
It is possible to perform ion exchange chromatography in bulk, on thin layers of medium such as glass or plastic plates coated with a layer of the desired stationary phase, or in chromatography columns. Thin layer chromatography or column chromatography share similarities in that they both act within the same governing principles; there is constant and frequent exchange of molecules as the mobile phase travels along the stationary phase. It is not imperative to add the sample in minute volumes as the predetermined conditions for the exchange column have been chosen so that there will be strong interaction between the mobile and stationary phases. Furthermore, the mechanism of the elution process will cause a compartmentalization of the differing molecules based on their respective chemical characteristics. This phenomenon is due to an increase in salt concentrations at or near the top of the column, thereby displacing the molecules at that position, while molecules bound lower are released at a later point when the higher salt concentration reaches that area. These principles are the reasons that ion exchange chromatography is an excellent candidate for initial chromatography steps in a complex purification procedure as it can quickly yield small volumes of target molecules regardless of a greater starting volume.
Comparatively simple devices are often used to apply counterions of increasing gradient to a chromatography column. Counterions such as copper (II) are chosen most often for effectively separating peptides and amino acids through complex formation.
A simple device can be used to create a salt gradient. Elution buffer is consistently being drawn from the chamber into the mixing chamber, thereby altering its buffer concentration. Generally, the buffer placed into the chamber is usually of high initial concentration, whereas the buffer placed into the stirred chamber is usually of low concentration. As the high concentration buffer from the left chamber is mixed and drawn into the column, the buffer concentration of the stirred column gradually increase. Altering the shapes of the stirred chamber, as well as of the limit buffer, allows for the production of concave, linear, or convex gradients of counterion.
A multitude of different mediums are used for the stationary phase. Among the most common immobilized charged groups used are trimethylaminoethyl (TAM), triethylaminoethyl (TEAE), diethyl-2-hydroxypropylaminoethyl (QAE), aminoethyl (AE), diethylaminoethyl (DEAE), sulpho (S), sulphomethyl (SM), sulphopropyl (SP), carboxy (C), and carboxymethyl (CM).
Successful packing of the column is an important aspect of ion chromatography. Stability and efficiency of a final column depends on packing methods, solvent used, and factors that affect mechanical properties of the column. In contrast to early inefficient dry- packing methods, wet slurry packing, in which particles that are suspended in an appropriate solvent are delivered into a column under pressure, shows significant improvement. Three different approaches can be employed in performing wet slurry packing: the balanced density method (solvent's density is about that of porous silica particles), the high viscosity method (a solvent of high viscosity is used), and the low viscosity slurry method (performed with low viscosity solvents).
Polystyrene is used as a medium for ion- exchange. It is made from the polymerization of styrene with the use of divinylbenzene and benzoyl peroxide. Such exchangers form hydrophobic interactions with proteins which can be irreversible. Due to this property, polystyrene ion exchangers are not suitable for protein separation. They are used on the other hand for the separation of small molecules in amino acid separation and removal of salt from water. Polystyrene ion exchangers with large pores can be used for the separation of protein but must be coated with a hydrophilic substance.
Cellulose based medium can be used for the separation of large molecules as they contain large pores. Protein binding in this medium is high and has low hydrophobic character. DEAE is an anion exchange matrix that is produced from a positive side group of diethylaminoethyl bound to cellulose or Sephadex.
Agarose gel based medium contain large pores as well but their substitution ability is lower in comparison to dextrans. The ability of the medium to swell in liquid is based on the cross-linking of these substances, the pH and the ion concentrations of the buffers used.
Incorporation of high temperature and pressure allows a significant increase in the efficiency of ion chromatography, along with a decrease in time. Temperature has an influence of selectivity due to its effects on retention properties. The retention factor (k = (tRg − tMg)/(tMg − text)) increases with temperature for small ions, and the opposite trend is observed for larger ions.
Despite ion selectivity in different mediums, further research is being done to perform ion exchange chromatography through the range of 40–175 °C.
An appropriate solvent can be chosen based on observations of how column particles behave in a solvent. Using an optical microscope, one can easily distinguish a desirable dispersed state of slurry from aggregated particles.
Weak and strong ion exchangers
A "strong" ion exchanger will not lose the charge on its matrix once the column is equilibrated and so a wide range of pH buffers can be used. "Weak" ion exchangers have a range of pH values in which they will maintain their charge. If the pH of the buffer used for a weak ion exchange column goes out of the capacity range of the matrix, the column will lose its charge distribution and the molecule of interest may be lost. Despite the smaller pH range of weak ion exchangers, they are often used over strong ion exchangers due to their having greater specificity. In some experiments, the retention times of weak ion exchangers are just long enough to obtain desired data at a high specificity.
Resins (often termed 'beads') of ion exchange columns may include functional groups such as weak/strong acids and weak/strong bases. There are also special columns that have resins with amphoteric functional groups that can exchange both cations and anions. Some examples of functional groups of strong ion exchange resins are quaternary ammonium cation (Q), which is an anion exchanger, and sulfonic acid (S, -SO2OH), which is a cation exchanger. These types of exchangers can maintain their charge density over a pH range of 0–14. Examples of functional groups of Weak ion exchange resins include diethylaminoethyl (DEAE, -C2H4N(C2H5)2), which is an anion exchanger, and carboxymethyl (CM, -CH2-COOH), which is a cation exchanger. These two types of exchangers can maintain the charge density of their columns over a pH range of 5–9.
In ion chromatography, the interaction of the solute ions and the stationary phase based on their charges determines which ions will bind and to what degree. When the stationary phase features positive groups which attracts anions, it is called an anion exchanger; when there are negative groups on the stationary phase, cations are attracted and it is a cation exchanger. The attraction between ions and stationary phase also depends on the resin, organic particles used as ion exchangers.
Each resin features relative selectivity which varies based on the solute ions present who will compete to bind to the resin group on the stationary phase. The selectivity coefficient, the equivalent to the equilibrium constant, is determined via a ratio of the concentrations between the resin and each ion, however, the general trend is that ion exchangers prefer binding to the ion with a higher charge, smaller hydrated radius, and higher polarizability, or the ability for the electron cloud of an ion to be disrupted by other charges. Despite this selectivity, excess amounts of an ion with a lower selectivity introduced to the column would cause the lesser ion to bind more to the stationary phase as the selectivity coefficient allows fluctuations in the binding reaction that takes place during ion exchange chromatography.
Following table shows the commonly used ion exchangers
Typical technique
A sample is introduced, either manually or with an autosampler, into a sample loop of known volume. A buffered aqueous solution known as the mobile phase carries the sample from the loop onto a column that contains some form of stationary phase material. This is typically a resin or gel matrix consisting of agarose or cellulose beads with covalently bonded charged functional groups. Equilibration of the stationary phase is needed in order to obtain the desired charge of the column. If the column is not properly equilibrated the desired molecule may not bind strongly to the column. The target analytes (anions or cations) are retained on the stationary phase but can be eluted by increasing the concentration of a similarly charged species that displaces the analyte ions from the stationary phase. For example, in cation exchange chromatography, the positively charged analyte can be displaced by adding positively charged sodium ions. The analytes of interest must then be detected by some means, typically by conductivity or UV/visible light absorbance.
Control an IC system usually requires a chromatography data system (CDS). In addition to IC systems, some of these CDSs can also control gas chromatography (GC) and HPLC.
Membrane exchange chromatography
A type of ion exchange chromatography, membrane exchange is a relatively new method of purification designed to overcome limitations of using columns packed with beads. Membrane Chromatographic devices are cheap to mass-produce and disposable unlike other chromatography devices that require maintenance and time to revalidate. There are three types of membrane absorbers that are typically used when separating substances. The three types are flat sheet, hollow fibre, and radial flow. The most common absorber and best suited for membrane chromatography is multiple flat sheets because it has more absorbent volume. It can be used to overcome mass transfer limitations and pressure drop, making it especially advantageous for isolating and purifying viruses, plasmid DNA, and other large macromolecules. The column is packed with microporous membranes with internal pores which contain adsorptive moieties that can bind the target protein. Adsorptive membranes are available in a variety of geometries and chemistry which allows them to be used for purification and also fractionation, concentration, and clarification in an efficiency that is 10 fold that of using beads. Membranes can be prepared through isolation of the membrane itself, where membranes are cut into squares and immobilized. A more recent method involved the use of live cells that are attached to a support membrane and are used for identification and clarification of signaling molecules.
Separating proteins
Ion exchange chromatography can be used to separate proteins because they contain charged functional groups. The ions of interest (in this case charged proteins) are exchanged for another ions (usually H+) on a charged solid support. The solutes are most commonly in a liquid phase, which tends to be water. Take for example proteins in water, which would be a liquid phase that is passed through a column. The column is commonly known as the solid phase since it is filled with porous synthetic particles that are of a particular charge. These porous particles are also referred to as beads, may be aminated (containing amino groups) or have metal ions in order to have a charge. The column can be prepared using porous polymers, for macromolecules of a mass of over 100 000 Da, the optimum size of the porous particle is about 1 μm2. This is because slow diffusion of the solutes within the pores does not restrict the separation quality. The beads containing positively charged groups, which attract the negatively charged proteins, are commonly referred to as anion exchange resins. The amino acids that have negatively charged side chains at pH 7 (pH of water) are glutamate and aspartate. The beads that are negatively charged are called cation exchange resins, as positively charged proteins will be attracted. The amino acids that have positively charged side chains at pH 7 are lysine, histidine and arginine.
The isoelectric point is the pH at which a compound - in this case a protein - has no net charge. A protein's isoelectric point or PI can be determined using the pKa of the side chains, if the amino (positive chain) is able to cancel out the carboxyl (negative) chain, the protein would be at its PI. Using buffers instead of water for proteins that do not have a charge at pH 7 is a good idea as it enables the manipulation of pH to alter ionic interactions between the proteins and the beads. Weakly acidic or basic side chains are able to have a charge if the pH is high or low enough respectively. Separation can be achieved based on the natural isoelectric point of the protein. Alternatively a peptide tag can be genetically added to the protein to give the protein an isoelectric point away from most natural proteins (e.g., 6 arginines for binding to a cation-exchange resin or 6 glutamates for binding to an anion-exchange resin such as DEAE-Sepharose).
Elution by increasing ionic strength of the mobile phase is more subtle. It works because ions from the mobile phase interact with the immobilized ions on the stationary phase, thus "shielding" the stationary phase from the protein, and letting the protein elute.
Elution from ion-exchange columns can be sensitive to changes of a single charge- chromatofocusing. Ion-exchange chromatography is also useful in the isolation of specific multimeric protein assemblies, allowing purification of specific complexes according to both the number and the position of charged peptide tags.
Gibbs–Donnan effect
In ion exchange chromatography, the Gibbs–Donnan effect is observed when the pH of the applied buffer and the ion exchanger differ, even up to one pH unit. For example, in anion-exchange columns, the ion exchangers repeal protons so the pH of the buffer near the column differs is higher than the rest of the solvent. As a result, an experimenter has to be careful that the protein(s) of interest is stable and properly charged in the "actual" pH.
This effect comes as a result of two similarly charged particles, one from the resin and one from the solution, failing to distribute properly between the two sides; there is a selective uptake of one ion over another. For example, in a sulphonated polystyrene resin, a cation exchange resin, the chlorine ion of a hydrochloric acid buffer should equilibrate into the resin. However, since the concentration of the sulphonic acid in the resin is high, the hydrogen of HCl has no tendency to enter the column. This, combined with the need of electroneutrality, leads to a minimum amount of hydrogen and chlorine entering the resin.
Uses
Clinical utility
A use of ion chromatography can be seen in argentation chromatography. Usually, silver and compounds containing acetylenic and ethylenic bonds have very weak interactions. This phenomenon has been widely tested on olefin compounds. The ion complexes the olefins make with silver ions are weak and made based on the overlapping of pi, sigma, and d orbitals and available electrons therefore cause no real changes in the double bond. This behavior was manipulated to separate lipids, mainly fatty acids from mixtures in to fractions with differing number of double bonds using silver ions. The ion resins were impregnated with silver ions, which were then exposed to various acids (silicic acid) to elute fatty acids of different characteristics.
Detection limits as low as 1 μM can be obtained for alkali metal ions.
It may be used for measurement of HbA1c, porphyrin and with water purification. Ion Exchange Resins(IER) have been widely used especially in medicines due to its high capacity and the uncomplicated system of the separation process. One of the synthetic uses is to use Ion Exchange Resins for kidney dialysis. This method is used to separate the blood elements by using the cellulose membraned artificial kidney.
Another clinical application of ion chromatography is in the rapid anion exchange chromatography technique used to separate creatine kinase (CK) isoenzymes from human serum and tissue sourced in autopsy material (mostly CK rich tissues were used such as cardiac muscle and brain). These isoenzymes include MM, MB, and BB, which all carry out the same function given different amino acid sequences. The functions of these isoenzymes are to convert creatine, using ATP, into phosphocreatine expelling ADP. Mini columns were filled with DEAE-Sephadex A-50 and further eluted with tris- buffer sodium chloride at various concentrations (each concentration was chosen advantageously to manipulate elution). Human tissue extract was inserted in columns for separation. All fractions were analyzed to see total CK activity and it was found that each source of CK isoenzymes had characteristic isoenzymes found within. Firstly, CK- MM was eluted, then CK-MB, followed by CK-BB. Therefore, the isoenzymes found in each sample could be used to identify the source, as they were tissue specific.
Using the information from results, correlation could be made about the diagnosis of patients and the kind of CK isoenzymes found in most abundant activity. From the finding, about 35 out of 71 patients studied suffered from heart attack (myocardial infarction) also contained an abundant amount of the CK-MM and CK-MB isoenzymes. Findings further show that many other diagnosis including renal failure, cerebrovascular disease, and pulmonary disease were only found to have the CK-MM isoenzyme and no other isoenzyme. The results from this study indicate correlations between various diseases and the CK isoenzymes found which confirms previous test results using various techniques. Studies about CK-MB found in heart attack victims have expanded since this study and application of ion chromatography.
Industrial applications
Since 1975 ion chromatography has been widely used in many branches of industry. The main beneficial advantages are reliability, very good accuracy and precision, high selectivity, high speed, high separation efficiency, and low cost of consumables. The most significant development related to ion chromatography are new sample preparation methods; improving the speed and selectivity of analytes separation; lowering of limits of detection and limits of quantification; extending the scope of applications; development of new standard methods; miniaturization and extending the scope of the analysis of a new group of substances. Allows for quantitative testing of electrolyte and proprietary additives of electroplating baths. It is an advancement of qualitative hull cell testing or less accurate UV testing. Ions, catalysts, brighteners and accelerators can be measured. Ion exchange chromatography has gradually become a widely known, universal technique for the detection of both anionic and cationic species. Applications for such purposes have been developed, or are under development, for a variety of fields of interest, and in particular, the pharmaceutical industry. The usage of ion exchange chromatography in pharmaceuticals has increased in recent years, and in 2006, a chapter on ion exchange chromatography was officially added to the United States Pharmacopia-National Formulary (USP-NF). Furthermore, in 2009 release of the USP-NF, the United States Pharmacopia made several analyses of ion chromatography available using two techniques: conductivity detection, as well as pulse amperometric detection. Majority of these applications are primarily used for measuring and analyzing residual limits in pharmaceuticals, including detecting the limits of oxalate, iodide, sulfate, sulfamate, phosphate, as well as various electrolytes including potassium, and sodium. In total, the 2009 edition of the USP-NF officially released twenty eight methods of detection for the analysis of active compounds, or components of active compounds, using either conductivity detection or pulse amperometric detection.
Drug development
There has been a growing interest in the application of IC in the analysis of pharmaceutical drugs. IC is used in different aspects of product development and quality control testing. For example, IC is used to improve stabilities and solubility properties of pharmaceutical active drugs molecules as well as used to detect systems that have higher tolerance for organic solvents. IC has been used for the determination of analytes as a part of a dissolution test. For instance, calcium dissolution tests have shown that other ions present in the medium can be well resolved among themselves and also from the calcium ion. Therefore, IC has been employed in drugs in the form of tablets and capsules in order to determine the amount of drug dissolve with time. IC is also widely used for detection and quantification of excipients or inactive ingredients used in pharmaceutical formulations. Detection of sugar and sugar alcohol in such formulations through IC has been done due to these polar groups getting resolved in ion column. IC methodology also established in analysis of impurities in drug substances and products. Impurities or any components that are not part of the drug chemical entity are evaluated and they give insights about the maximum and minimum amounts of drug that should be administered in a patient per day.
| Physical sciences | Chromatography | Chemistry |
4576465 | https://en.wikipedia.org/wiki/Flower | Flower | A flower, also known as a bloom or blossom, is the reproductive structure found in flowering plants (plants of the division Angiospermae). Flowers consist of a combination of vegetative organs – sepals that enclose and protect the developing flower. Petals attract pollinators, and reproductive organs that produce gametophytes, which in flowering plants produce gametes. The male gametophytes, which produce sperm, are enclosed within pollen grains produced in the anthers. The female gametophytes are contained within the ovules produced in the ovary. In some plants, multiple flowers occur singly on a pedicel (flower stalk), and some are arranged in a group (inflorescence) on a peduncle (inflorescence stalk).
Most flowering plants depend on animals, such as bees, moths, and butterflies, to transfer their pollen between different flowers, and have evolved to attract these pollinators by various strategies, including brightly colored, large petals with patterns only visible to under ultraviolet light, attractive scents, and the production of nectar, a food source for pollinators. In this way, many flowering plants have co-evolved with pollinators to be mutually dependent on services they provide to one another—in the plant's case, a means of reproduction; in the pollinator's case, a source of food.
When pollen from the anther of a flower is transferred to the stigma to another, it is called pollination. Some flowers may self-pollinate, producing seed using pollen from a different flower of the same plant, but others have mechanisms to prevent self-pollination and rely on cross-pollination, when pollen is transferred from the anther of one flower to the stigma of another flower on a different individual of the same species. Self-pollination happens in flowers where the stamen and carpel mature at the same time, and are positioned so that the pollen can land on the flower's stigma. This pollination does not require an investment from the plant to provide nectar and pollen as food for pollinators. Some flowers produce diaspores without fertilization (parthenocarpy). After fertilization, the ovary of the flower develops into fruit containing seeds.
Flowers have long been appreciated for their beauty and pleasant scents, and also hold cultural significance as religious, ritual, or symbolic objects, or sources of medicine and food.
Etymology
Flower is from the Middle English , which referred to both the ground grain and the reproductive structure in plants, before splitting off in the 17th century. It comes originally from the Latin name of the Italian goddess of flowers, Flora. The early word for flower in English was blossom, though it now refers to flowers only of fruit trees.
Morphology
The morphology of a flower, or its form and structure, can be considered in two parts: the vegetative part, consisting of non-reproductive structures such as petals; and the reproductive or sexual parts. A stereotypical flower is made up of four kinds of structures arranged in whorls around the tip of a short stalk or axis, called a receptacle. The four main whorls (starting from the base of the flower or lowest node and working upwards) are the calyx, corolla, androecium, and gynoecium. Together the calyx and corolla make up the non-reproductive part of the flower called the perianth, and in monocotyledons, may not be differentiated. If this is the case, then they are described as tepals.
Perianth
Calyx
The sepals, collectively called the calyx, are modified leaves that occur on the outermost whorl of the flower. They are leaf-like, in that they have a broad base, stomata and chlorophyll and may have stipules. Sepals are often waxy and tough, and grow quickly to protect the flower as it develops. They may be deciduous, but will more commonly grow on to assist in fruit dispersal. If the calyx is fused it is called gamosepalous.
Corolla
The petals, collectively called corolla, are almost or completely fiberless leaf-like structures that form the innermost whorl of the perianth. They are often delicate and thin and are usually colored, shaped, or scented to encourage pollination. Although similar to leaves in shape, they are more comparable to stamens in that they form almost simultaneously with one another, but their subsequent growth is delayed. If the corolla is fused together it is called sympetalous. In monocotyledonous flowers (e.g., Lilium sp.), petals and sepals are indistinguishable and are individually called tepals. Petals also tend to have patterns only visible under ultraviolet light, which are visible to pollinators but not to humans.
Reproductive
Androecium
The androecium, consisting of stamens, is the whorl of pollen-producing male parts. Stamens consist typically of an anther, made up of four pollen sacs arranged in two thecae, connected to a filament, or stalk. The anther contains microsporocytes which become pollen, the male gametophyte, after undergoing meiosis. Although they exhibit the widest variation among floral organs, the androecium is usually confined just to one whorl and to two whorls only in rare cases. Stamens range in number, size, shape, orientation, and in their point of connection to the flower.
In general, there is only one type of stamen, but there are plant species where the flowers have two types; a "normal" one and one with anthers that produce sterile pollen meant to attract pollinators.
Gynoecium
The gynoecium, consisting of one or more carpels, is the female part of the flower found on the innermost whorl. Each carpel consists of a stigma, which receives pollen, a style, which acts as a stalk, and an ovary, which contains the ovules. Carpels may occur in one to several whorls, and when fused are often described as a pistil. Inside the ovary, the ovules are attached to the placenta by structures called funiculi.
Variation
Although this arrangement is considered "typical", plant species show a wide variation in floral structure. The four main parts of a flower are generally defined by their positions on the receptacle and not by their function. Many flowers lack some parts or parts may be modified into other functions or look like what is typically another part. In some families, such as the grasses, the petals are greatly reduced; in many species, the sepals are colorful and petal-like. Other flowers have modified petal-like stamens; the double flowers of peonies and roses are mostly petaloid stamens.
Many flowers have symmetry. When the perianth is bisected through the central axis from any point and symmetrical halves are produced, the flower is said to be actinomorphic or regular. This is an example of radial symmetry. When flowers are bisected and produce only one line that produces symmetrical halves, the flower is said to be irregular or zygomorphic. If, in rare cases, they have no symmetry at all they are called asymmetric.
Flowers may be directly attached to the plant at their base (sessile—the supporting stalk or stem is highly reduced or absent). The stem or stalk subtending a flower, or an inflorescence of flowers, is called a peduncle. If a peduncle supports more than one flower, the stems connecting each flower to the main axis are called pedicels. The apex of a flowering stem forms a terminal swelling which is called the torus or receptacle.
In the majority of species, individual flowers have both carpels and stamens. These flowers are described by botanists as being perfect, bisexual, or hermaphrodite. In some species of plants, the flowers are imperfect or unisexual: having only either male (stamen) or female (carpel) parts. If unisexual male and female flowers appear on the same plant, the species is called monoecious. However, if an individual plant is either female or male, the species is called dioecious. Many flowers have floral nectaries, which are glands that produce a sugary fluid (nectar) used to attract pollinators. They are not considered as an organ on their own.
Inflorescence
In those species that have more than one flower on an axis, the collective cluster of flowers is called an inflorescence. Some inflorescences are composed of many small flowers arranged in a formation that resembles a single flower. A common example of this is most members of the very large composite (Asteraceae) group. A single daisy or sunflower, for example, is not a flower but a flower head—an inflorescence composed of numerous flowers (or florets). An inflorescence may include specialized stems and modified leaves known as bracts.
Floral diagrams and formulae
A floral formula is a way to represent the structure of a flower using specific letters, numbers, and symbols, presenting substantial information about the flower in a compact form. It can represent a taxon, usually giving ranges of the numbers of different organs, or particular species. Floral formulae have been developed in the early 19th century and their use has declined since. Prenner et al. (2010) devised an extension of the existing model to broaden the descriptive capability of the formula. The format of floral formulae differs in different parts of the world, yet they convey the same information.
The structure of a flower can also be expressed by the means of floral diagrams. The use of schematic diagrams can replace long descriptions or complicated drawings as a tool for understanding both floral structure and evolution. Such diagrams may show important features of flowers, including the relative positions of the various organs, including the presence of fusion and symmetry, as well as structural details.
Development
A flower develops on a modified shoot or axis from a determinate apical meristem (determinate meaning the axis grows to a set size). It has compressed internodes, bearing structures that in classical plant morphology are interpreted as highly modified leaves. Detailed developmental studies, however, have shown that stamens are often initiated more or less like modified stems (caulomes) that in some cases may even resemble branchlets. Taking into account the whole diversity in the development of the androecium of flowering plants, we find a continuum between modified leaves (phyllomes), modified stems (caulomes), and modified branchlets (shoots).
Transition
The transition to flowering is one of the major phase changes that a plant makes during its life cycle. The transition must take place at a time that is favorable for fertilization and the formation of seeds, hence ensuring maximal reproductive success. To meet these needs a plant can interpret important endogenous and environmental cues such as changes in levels of plant hormones and seasonable temperature and photoperiod changes. Many perennial and most biennial plants require vernalization to flower. The molecular interpretation of these signals is through the transmission of a complex signal known as florigen, which involves a variety of genes, including Constans, Flowering Locus C, and Flowering Locus T. Florigen is produced in the leaves in reproductively favorable conditions and acts in buds and growing tips to induce several different physiological and morphological changes.The first step of the transition is the transformation of the vegetative stem primordia into floral primordia. This occurs as biochemical changes take place to change the cellular differentiation of leaf, bud and stem tissues into tissue that will grow into the reproductive organs. Growth of the central part of the stem tip stops or flattens out and the sides develop protuberances in a whorled or spiral fashion around the outside of the stem end. These protuberances develop into the sepals, petals, stamens, and carpels. Once this process begins, in most plants, it cannot be reversed and the stems develop flowers, even if the initial start of the flower formation event was dependent on some environmental cue.
Organ development
The ABC model is a simple model that describes the genes responsible for the development of flowers. Three gene activities interact in a combinatorial manner to determine the developmental identities of the primordia organ within the floral apical meristem. These gene functions are called A, B, and C. Genes are expressed in only the outer and lower most section of the apical meristem, which becomes a whorl of sepals. In the second whorl, both A and B genes are expressed, leading to the formation of petals. In the third whorl, B and C genes interact to form stamens and in the center of the flower C genes alone give rise to carpels. The model is based upon studies of aberrant flowers and mutations in Arabidopsis thaliana and the snapdragon, Antirrhinum majus. For example, when there is a loss of B gene function, mutant flowers are produced with sepals in the first whorl as usual, but also in the second whorl instead of the normal petal formation. In the third whorl, the lack of the B function but the presence of the C function mimics the fourth whorl, leading to the formation of carpels also in the third whorl.
Function
The principal purpose of a flower is the reproduction of the individual and the species. All flowering plants are heterosporous, that is, every individual plant produces two types of spores. Microspores are produced by meiosis inside anthers and megaspores are produced inside ovules that are within an ovary. Anthers typically consist of four microsporangia and an ovule is an integumented megasporangium. Both types of spores develop into gametophytes inside sporangia. As with all heterosporous plants, the gametophytes also develop inside the spores, i.e., they are endosporic.
Pollination
Since the flowers are the reproductive organs of the plant, they mediate the joining of the sperm, contained within pollen, to the ovules — contained in the ovary. Pollination is the movement of pollen from the anthers to the stigma. Normally pollen is moved from one plant to another, known as cross-pollination, but many plants can self-pollinate. Cross-pollination is preferred because it allows for genetic variation, which contributes to the survival of the species. Many flowers depend on external factors for pollination, such as the wind, water, animals, and especially insects. Larger animals such as birds, bats, and even some pygmy possums, however, can also be employed. To accomplish this, flowers have specific designs which encourage the transfer of pollen from one plant to another of the same species. The period during which this process can take place (when the flower is fully expanded and functional) is called anthesis, hence the study of pollination biology is called anthecology.
Flowering plants usually face evolutionary pressure to optimize the transfer of their pollen, and this is typically reflected in the morphology of the flowers and the behavior of the plants. Pollen may be transferred between plants via several 'vectors,' or methods. Around 80% of flowering plants make use of biotic or living vectors. Others use abiotic, or non-living, vectors and some plants make use of multiple vectors, but most are highly specialized.
Though some fit between or outside of these groups, most flowers can be divided between the following two broad groups of pollination methods:
Biotic pollination
Flowers that use biotic vectors attract and use insects, bats, birds, or other animals to transfer pollen from one flower to the next. Often they are specialized in shape and have an arrangement of the stamens that ensures that pollen grains are transferred to the bodies of the pollinator when it lands in search of its attractant (such as nectar, pollen, or a mate). In pursuing this attractant from many flowers of the same species, the pollinator transfers pollen to the stigmas—arranged with equally pointed precision—of all of the flowers it visits. Many flowers rely on simple proximity between flower parts to ensure pollination, while others have elaborate designs to ensure pollination and prevent self-pollination. Flowers use animals including: insects (entomophily), birds (ornithophily), bats (chiropterophily), lizards, and even snails and slugs (malacophilae).
Attraction methods
Plants cannot move from one location to another, thus many flowers have evolved to attract animals to transfer pollen between individuals in dispersed populations. Most commonly, flowers are insect-pollinated, known as entomophilous; literally "insect-loving" in Greek. To attract these insects flowers commonly have glands called nectaries on various parts that attract animals looking for nutritious nectar. Some flowers have glands called elaiophores, which produce oils rather than nectar. Birds and bees have color vision, enabling them to seek out colorful flowers. Some flowers have patterns, called nectar guides, that show pollinators where to look for nectar; they may be visible only under ultraviolet light, which is visible to bees and some other insects.
Flowers also attract pollinators by scent, though not all flower scents are appealing to humans; several flowers are pollinated by insects that are attracted to rotten flesh and have flowers that smell like dead animals. These are often called carrion flowers, including plants in the genus Rafflesia, and the titan arum. Flowers pollinated by night visitors, including bats and moths, are likely to concentrate on scent to attract pollinators and so most such flowers are white. Some plants pollinated by bats have a sonar-reflecting petal above its flowers, which helps the bat find them, and one species, the cactus Espostoa frutescens, has flowers that are surrounded by an area of sound-absorbent and woolly hairs called the cephalium, which absorbs the bat's ultrasound instead.
Flowers are also specialized in shape and have an arrangement of the stamens that ensures that pollen grains are transferred to the bodies of the pollinator when it lands in search of its attractant. Other flowers use mimicry or pseudocopulation to attract pollinators. Many orchids, for example, produce flowers resembling female bees or wasps in color, shape, and scent. Males move from one flower to the next in search of a mate, pollinating the flowers.
Pollinator relationships
Many flowers have close relationships with one or a few specific pollinating organisms. Many flowers, for example, attract only one specific species of insect and therefore rely on that insect for successful reproduction. This close relationship is an example of coevolution, as the flower and pollinator have developed together over a long period to match each other's needs. This close relationship compounds the negative effects of extinction, however, since the extinction of either member in such a relationship would almost certainly mean the extinction of the other member as well.
Abiotic pollination
Flowers that use abiotic, or non-living, vectors use the wind or, much less commonly, water, to move pollen from one flower to the next. In wind-dispersed (anemophilous) species, the tiny pollen grains are carried, sometimes many thousands of kilometers, by the wind to other flowers. Common examples include the grasses, birch trees, along with many other species in the order Fagales, ragweeds, and many sedges. They do not need to attract pollinators and therefore tend not to grow large, showy, or colorful flowers, and do not have nectaries, nor a noticeable scent. Because of this, plants typically have many thousands of tiny flowers which have comparatively large, feathery stigmas; to increase the chance of pollen being received. Whereas the pollen of entomophilous flowers is usually large, sticky, and rich in protein (to act as a "reward" for pollinators), anemophilous flower pollen is typically small-grained, very light, smooth, and of little nutritional value to insects. In order for the wind to effectively pick up and transport the pollen, the flowers typically have anthers loosely attached to the end of long thin filaments, or pollen forms around a catkin which moves in the wind. Rarer forms of this involve individual flowers being moveable by the wind (pendulous), or even less commonly; the anthers exploding to release the pollen into the wind.
Pollination through water (hydrophily) is a much rarer method, occurring in only around 2% of abiotically pollinated flowers. Common examples of this include Calitriche autumnalis, Vallisneria spiralis and some sea-grasses. One characteristic which most species in this group share is a lack of an exine, or protective layer, around the pollen grain. Paul Knuth identified two types of hydrophilous pollination in 1906 and Ernst Schwarzenbach added a third in 1944. Knuth named his two groups 'Hyphydrogamy' and the more common 'Ephydrogamy'. In hyphydrogamy pollination occurs below the surface of the water and so the pollen grains are typically negatively buoyant. For marine plants that exhibit this method, the stigmas are usually stiff, while freshwater species have small and feathery stigmas. In ephydrogamy pollination occurs on the surface of the water and so the pollen has a low density to enable floating, though many also use rafts, and are hydrophobic. Marine flowers have floating thread-like stigmas and may have adaptations for the tide, while freshwater species create indentations in the water. The third category, set out by Schwarzenbach, is those flowers which transport pollen above the water through conveyance. This ranges from floating plants, (Lemnoideae), to staminate flowers (Vallisneria). Most species in this group have dry, spherical pollen which sometimes forms into larger masses, and female flowers which form depressions in the water; the method of transport varies.
Mechanisms
Flowers can be pollinated by two mechanisms; cross-pollination and self-pollination. No mechanism is indisputably better than the other as they each have their advantages and disadvantages. Plants use one or both of these mechanisms depending on their habitat and ecological niche.
Cross-pollination
Cross-pollination is the pollination of the carpel by pollen from a different plant of the same species. Because the genetic make-up of the sperm contained within the pollen from the other plant is different, their combination will result in a new, genetically distinct, plant, through the process of sexual reproduction. Since each new plant is genetically distinct, the different plants show variation in their physiological and structural adaptations and so the population as a whole is better prepared for an adverse occurrence in the environment. Cross-pollination, therefore, increases the survival of the species and is usually preferred by flowers for this reason.
The principal adaptive function of flowers is the promotion of cross-pollination or outcrossing, a process that allows the masking of deleterious mutations in the genome of progeny. The masking effect of outcrossing sexual reproduction is known as "genetic complementation". This beneficial effect of outcrossing on progeny is also recognized as hybrid vigour or heterosis. Once outcrossing is established due to the benefits of genetic complementation, subsequent switching to inbreeding becomes disadvantageous because it allows the expression of the previously masked deleterious recessive mutations, usually referred to as inbreeding depression. Charles Darwin in his 1889 book The Effects of Cross and Self-Fertilization in the Vegetable Kingdom at the beginning of chapter XII noted, "The first and most important of the conclusions which may be drawn from the observations given in this volume, is that generally cross-fertilisation is beneficial and self-fertilisation often injurious, at least with the plants on which I experimented."
Self-pollination
Self-pollination is the pollination of the carpel of a flower by pollen from either the same flower or another flower on the same plant, leading to the creation of a genetic clone through asexual reproduction. This increases the reliability of producing seeds, the rate at which they can be produced, and lowers the amount energy needed. But, most importantly, it limits genetic variation. In addition, self-pollination causes inbreeding depression, due largely to the expression of recessive deleterious mutations.
The extreme case of self-fertilization, when the ovule is fertilized by pollen from the same flower or plant, occurs in flowers that always self-fertilize, such as many dandelions. Some flowers are self-pollinated and have flowers that never open or are self-pollinated before the flowers open; these flowers are called cleistogamous; many species in the genus Viola exhibit this, for example.
Conversely, many species of plants have ways of preventing self-pollination and hence, self-fertilization. Unisexual male and female flowers on the same plant may not appear or mature at the same time, or pollen from the same plant may be incapable of fertilizing its ovules. The latter flower types, which have chemical barriers to their own pollen, are referred to as self-incompatible. In Clianthus puniceus, self-pollination is used strategically as an "insurance policy". When a pollinator, in this case a bird, visits C. puniceus, it rubs off the stigmatic covering and allows for pollen from the bird to enter the stigma. If no pollinators visit, however, then the stigmatic covering falls off naturally to allow for the flower's own anthers to pollinate the flower through self-pollination.
Allergies
Pollen is a large contributor to asthma and other respiratory allergies which combined affect between 10 and 50% of people worldwide. This number appears to be growing, as the temperature increases due to climate change mean that plants are producing more pollen, which is also more allergenic. Pollen is difficult to avoid, however, because of its small size and prevalence in the natural environment. Most of the pollen which causes allergies is that produced by wind-dispersed pollinators such as the grasses, birch trees, oak trees, and ragweeds; the allergens in pollen are proteins which are thought to be necessary in the process of pollination.
Fertilization
Fertilization, also called Syngamy, is preceded by pollination, which is the movement of pollen from the stamen to the carpel. It encompasses both plasmogamy, the fusion of the protoplasts, and karyogamy, the fusion of the nuclei. When pollen lands on the stigma of the flower it begins creating a pollen tube which runs down through the style and into the ovary. After penetrating the center-most part of the ovary it enters the egg apparatus and into one synergid. At this point the end of the pollen tube bursts and releases the two sperm cells, one of which makes its way to an egg, while also losing its cell membrane and much of its protoplasm. The sperm's nucleus then fuses with the egg's nucleus, resulting in the formation of a zygote, a diploid (two copies of each chromosome) cell.
In Angiosperms (flowering plants) a process known as double fertilization, which involves both karyogamy and plasmogamy, occurs. In double fertilization the second sperm cell subsequently also enters the synergid and fuses with the two polar nuclei of the central cell. Since all three nuclei are haploid, they result in a large endosperm nucleus which is triploid.
Seed development
Following the formation of zygote it begins to grow through nuclear and cellular divisions, called mitosis, eventually becoming a small group of cells. One section of it becomes the embryo, while the other becomes the suspensor; a structure which forces the embryo into the endosperm and is later undetectable. Two small primordia also form at this time, that later become the cotyledon, which is used as an energy store. Plants which grow out one of these primordia are called monocotyledons, while those that grow out two are dicotyledons. The next stage is called the Torpedo stage and involves the growth of several key structures, including: the radicle (embryotic root), the epicotyl (embryotic stem), and the hypocotyl, (the root/shoot junction). In the final step vascular tissue develops around the seed.
Fruit development
The ovary, inside which the seed is forming from the ovule, grows into a fruit. All the other main floral parts wither and die during this development, including: the style, stigma, sepals, stamens, and petals. The fruit contains three structures: the exocarp, or outer layer, the mesocarp, or the fleshy part, and the endocarp, or innermost layer, while the fruit wall is called the pericarp. The size, shape, toughness and thickness varies among different dry and fleshy fruits. This is because it is directly connected to the method of seed dispersal; that being the purpose of fruit - to encourage or enable the seed's dispersal and protect the seed while doing so.
Seed dispersal
Following the pollination of a flower, fertilization, and finally the development of a seed and fruit, a mechanism is typically used to disperse the fruit away from the plant. In Angiosperms (flowering plants) seeds are dispersed away from the plant so as to not force competition between the mother and the daughter plants, as well as to enable the colonization of new areas. They are often divided into two categories, though many plants fall in between or in one or more of these:
Allochory
In allochory, plants use an external vector, or carrier, to transport their seeds away from them. These can be either biotic (living), such as by birds and ants, or abiotic (non-living), such as by the wind or water.
Biotic vectors
Many plants use biotic vectors to disperse their seeds away from them. This method falls under the umbrella term zoochory, while endozoochory, also known as fruigivory, refers specifically to plants adapted to grow fruit in order to attract animals to eat them. Once eaten they go through typically go through animal's digestive system and are dispersed away from the plant. Some seeds are specially adapted either to last in the gizzard of animals or even to germinate better after passing through them. They can be eaten by birds (ornithochory), bats (chiropterochory), rodents, primates, ants (myrmecochory), non-bird sauropsids (saurochory), mammals in general (mammaliochory), and even fish. Typically their fruit are fleshy, have a high nutritional value, and may have chemical attractants as an additional "reward" for dispersers. This is reflected morphologically in the presence of more pulp, an aril, and sometimes an elaiosome (primarily for ants), which are other fleshy structures.
Epizoochory occurs in plants whose seeds are adapted to cling on to animals and be dispersed that way, such as many species in the genus Acaena. Typically these plants seed's have hooks or a viscous surface to easier grip to animals, which include birds and animals with fur. Some plants use mimesis, or imitation, to trick animals into dispersing the seeds and these often have specially adapted colors.
The final type of zoochory is called synzoochory, which involves neither the digestion of the seeds, nor the unintentional carrying of the seed on the body, but the deliberate carrying of the seeds by the animals. This is usually in the mouth or beak of the animal (called Stomatochory), which is what is used for many birds and all ants.
Abiotic vectors
In abiotic dispersal plants use the vectors of the wind, water, or a mechanism of their own to transport their seeds away from them. Anemochory involves using the wind as a vector to disperse plant's seeds. Because these seeds have to travel in the wind, they are almost always small — sometimes even dust-like, have a high surface-area-to-volume ratio, and are produced in a large number — sometimes up to a million. Plants such as tumbleweeds detach the entire shoot to let the seeds roll away with the wind. Another common adaptation are wings, plumes or balloon-like structures that let the seeds stay in the air for longer and hence travel farther.
In hydrochory plants are adapted to disperse their seeds through bodies of water and so typically are buoyant and have a low relative density with regards to the water. Commonly seeds are adapted morphologically with hydrophobic surfaces, small size, hairs, slime, oil, and sometimes air spaces within the seeds. These plants fall into three categories: ones where seeds are dispersed on the surface of water currents, under the surface of water currents, and by rain landing on a plant.
Autochory
In autochory, plants create their own vectors to transport the seeds away from them. Adaptations for this usually involve the fruits exploding and forcing the seeds away ballistically, such as in Hura crepitans, or sometimes in the creation of creeping diaspores. Because of the relatively small distances that these methods can disperse their seeds, they are often paired with an external vector.
Evolution
While land plants have existed for about 425 million years, the first ones reproduced by a simple adaptation of their aquatic counterparts: spores. In the sea, plants—and some animals—can simply scatter out genetic clones of themselves to float away and grow elsewhere. This is how early plants reproduced. But plants soon evolved methods of protecting these copies to deal with drying out and other damage which is even more likely on land than in the sea. The protection became the seed, though it had not yet evolved the flower. Early seed-bearing plants include the ginkgo and conifers.
Several groups of extinct gymnosperms, particularly seed ferns, have been proposed as the ancestors of flowering plants but there is no continuous fossil evidence showing exactly how flowers evolved. The apparently sudden appearance of relatively modern flowers in the fossil record posed such a problem for the theory of evolution that it was called an "abominable mystery" by Charles Darwin.
Recently discovered angiosperm fossils such as Archaefructus, along with further discoveries of fossil gymnosperms, suggest how angiosperm characteristics may have been acquired in a series of steps. An early fossil of a flowering plant, Archaefructus liaoningensis from China, is dated about 125 million years old. Even earlier from China is the 125–130 million years old Archaefructus sinensis. In 2015 a plant (130 million-year-old Montsechia vidalii, discovered in Spain) was claimed to be 130 million years old. In 2018, scientists reported that the earliest flowers began about 180 million years ago.
Recent DNA analysis (molecular systematics) shows that Amborella trichopoda, found on the Pacific island of New Caledonia, is the only species in the sister group to the rest of the flowering plants, and morphological studies suggest that it has features which may have been characteristic of the earliest flowering plants.
Besides the hard proof of flowers in or shortly before the Cretaceous, there is some circumstantial evidence of flowers as much as 250 million years ago. A chemical used by plants to defend their flowers, oleanane, has been detected in fossil plants that old, including gigantopterids, which evolved at that time and bear many of the traits of modern, flowering plants, though they are not known to be flowering plants themselves, because only their stems and prickles have been found preserved in detail; one of the earliest examples of petrification.
The similarity in leaf and stem structure can be very important, because flowers are genetically just an adaptation of normal leaf and stem components on plants, a combination of genes normally responsible for forming new shoots. The most primitive flowers are thought to have had a variable number of flower parts, often separate from (but in contact with) each other. The flowers would have tended to grow in a spiral pattern, to be bisexual (in plants, this means both male and female parts on the same flower), and to be dominated by the ovary (female part). As flowers grew more advanced, some variations developed parts fused together, with a much more specific number and design, and with either specific sexes per flower or plant, or at least "ovary inferior".
The general assumption is that the function of flowers, from the start, was to involve animals in the reproduction process. Pollen can be scattered without bright colors and obvious shapes, which would therefore be a liability, using the plant's resources, unless they provide some other benefit. One proposed reason for the sudden, fully developed appearance of flowers is that they evolved in an isolated setting like an island, or chain of islands, where the plants bearing them were able to develop a highly specialized relationship with some specific animal (a wasp, for example), the way many island species develop today. This symbiotic relationship, with a hypothetical wasp bearing pollen from one plant to another much the way fig wasps do today, could have eventually resulted in both the plant(s) and their partners developing a high degree of specialization. Island genetics is believed to be a common source of speciation, especially when it comes to radical adaptations which seem to have required inferior transitional forms. Note that the wasp example is not incidental; bees, apparently evolved specifically for symbiotic plant relationships, are descended from wasps.
Likewise, most fruit used in plant reproduction comes from the enlargement of parts of the flower. This fruit is frequently a tool which depends upon animals wishing to eat it, and thus scattering the seeds it contains.
While many such symbiotic relationships remain too fragile to survive competition with mainland organisms, flowers proved to be an unusually effective means of production, spreading (whatever their actual origin) to become the dominant form of land plant life.
Flower evolution continues to the present day; modern flowers have been so profoundly influenced by humans that many of them cannot be pollinated in nature. Many modern, domesticated flowers used to be simple weeds, which only sprouted when the ground was disturbed. Some of them tended to grow with human crops, and the prettiest did not get plucked because of their beauty, developing a dependence upon and special adaptation to human affection.
Colour
Many flowering plants reflect as much light as possible within the range of visible wavelengths of the pollinator the plant intends to attract. Flowers that reflect the full range of visible light are generally perceived as white by a human observer. An important feature of white flowers is that they reflect equally across the visible spectrum. While many flowering plants use white to attract pollinators, the use of color is also widespread (even within the same species). Color allows a flowering plant to be more specific about the pollinator it seeks to attract. The color model used by human color reproduction technology (CMYK) relies on the modulation of pigments that divide the spectrum into broad areas of absorption. Flowering plants by contrast are able to shift the transition point wavelength between absorption and reflection. If it is assumed that the visual systems of most pollinators view the visible spectrum as circular then it may be said that flowering plants produce color by absorbing the light in one region of the spectrum and reflecting the light in the other region. With CMYK, color is produced as a function of the amplitude of the broad regions of absorption. Flowering plants by contrast produce color by modifying the frequency (or rather wavelength) of the light reflected. Most flowers absorb light in the blue to yellow region of the spectrum and reflect light from the green to red region of the spectrum. For many species of flowering plant, it is the transition point that characterizes the color that they produce. Color may be modulated by shifting the transition point between absorption and reflection and in this way a flowering plant may specify which pollinator it seeks to attract. Some flowering plants also have a limited ability to modulate areas of absorption. This is typically not as precise as control over wavelength. Humans observers will perceive this as degrees of saturation (the amount of white in the color).
Classical taxonomy
In plant taxonomy, which is the study of the classification and identification of plants, the morphology of plant's flowers are used extensively – and have been for thousands of years. Although the history of plant taxonomy extends back to at least around 300 B.C. with the writings of Theophrastus, the foundation of the modern science is based on works in the 18th and 19th centuries.
Carl Linnaeus (1707–1778), was a Swedish botanist who spent most of his working life as a professor of natural history. His landmark 1757 book Species Plantarum lays out his system of classification as well as the concept of binomial nomenclature, the latter of which is still used around the world today. He identified 24 classes, based mainly on the number, length and union of the stamens.
The first ten classes follow the number of stamens directly (Octandria have 8 stamens etc.), while class eleven has 11–20 stamens and classes twelve and thirteen have 20 stamens; differing only in their point of attachment. The next five classes deal with the length of the stamens and the final five with the nature of the reproductive capability of the plant; where the stamen grows; and if the flower is concealed or exists at all (such as in ferns). This method of classification, despite being artificial, was used extensively for the following seven decades, before being replaced by the system of another botanist.
Antoine Laurent de Jussieu (1748–1836) was a French botanist whose 1787 work Genera plantarum: secundum ordines naturales disposita set out a new method for classifying plants; based instead on natural characteristics. Plants were divided by the number, if any, of cotyledons, and the location of the stamens.
The next most major system of classification came in the late 19th century from the botanists Joseph Dalton Hooker (1817–1911) and George Bentham (1800–1884). They built on the earlier works of de Jussieu and Augustin Pyramus de Candolle and devised a system which is still used in many of the world's herbaria.
Plants were divided at the highest level by the number of cotyledons and the nature of the flowers, before falling into orders (families), genera, and species. This system of classification was published in their Genera plantarum in three volumes between 1862 and 1883. It is the most highly regarded and deemed the "best system of classification," in some settings.
Following the development in scientific thought after Darwin's On the Origin of Species, many botanists have used more phylogenetic methods and the use of genetic sequencing, cytology, and palynology has become increasingly common. Despite this, morphological characteristics such as the nature of the flower and inflorescence still make up the bedrock of plant taxonomy.
Symbolism
Many flowers have important symbolic meanings in Western culture. The practice of assigning meanings to flowers is known as floriography. Some of the more common examples include:
Red roses are given as a symbol of love, beauty, and passion.
Poppies are a symbol of consolation in time of death. In the United Kingdom, New Zealand, Australia and Canada, red poppies are worn to commemorate soldiers who have died in times of war.
Irises/Lily are used in burials as a symbol referring to "resurrection/life". It is also associated with stars (sun) and its petals blooming/shining.
Daisies are a symbol of innocence.
Because of their varied and colorful appearance, flowers have long been a favorite subject of visual artists as well. Some of the most celebrated paintings from well-known painters are of flowers, such as Van Gogh's sunflowers series or Monet's water lilies. Flowers are also dried, freeze dried and pressed in order to create permanent, three-dimensional pieces of floral art.
Flowers within art are also representative of the female genitalia, as seen in the works of artists such as Georgia O'Keeffe, Imogen Cunningham, Veronica Ruiz de Velasco, and Judy Chicago, and in fact in Asian and western classical art. Many cultures around the world have a marked tendency to associate flowers with femininity.
The great variety of delicate and beautiful flowers has inspired the works of numerous poets, especially from the 18th–19th century Romantic era. Famous examples include William Wordsworth's I Wandered Lonely as a Cloud and William Blake's Ah! Sun-Flower.
Their symbolism in dreams has also been discussed, with possible interpretations including "blossoming potential".
The Roman goddess of flowers, gardens, and the season of Spring is Flora. The Greek goddess of spring, flowers and nature is Chloris.
In Hindu mythology, flowers have a significant status. Vishnu, one of the three major gods in the Hindu system, is often depicted standing straight on a lotus flower. Apart from the association with Vishnu, the Hindu tradition also considers the lotus to have spiritual significance. For example, it figures in the Hindu stories of creation.
Human use
History shows that flowers have been used by humans for thousands of years, to serve a variety of purposes. An early example of this is from about 4,500 years ago in Ancient Egypt, where flowers would be used to decorate women's hair. Flowers have also inspired art time and time again, such as in Monet's Water Lilies or William Wordsworth's poem about daffodils entitled: "I Wandered Lonely as a Cloud".
In modern times, people have sought ways to cultivate, buy, wear, or otherwise be around flowers and blooming plants, partly because of their agreeable appearance and smell. Around the world, people use flowers to mark important events in their lives:
For new births or christenings
As a corsage or boutonniere worn at social functions or for holidays
As tokens of love or esteem
For wedding flowers for the bridal party, and as decorations for wedding venues
As brightening decorations within the home
As a gift of remembrance for bon voyage parties, welcome-home parties, and "thinking of you" gifts
For funeral flowers and expressions of sympathy for the grieving
For worship. In Christianity, chancel flowers often adorn churches. In Hindu culture, adherents commonly bring flowers as a gift to temples
Flowers like jasmine have been used as a replacement for traditional tea in China for centuries. Most recently many other herbs and flowers used traditionally across the world are gaining importance to preapare a range of floral tea.
People therefore grow flowers around their homes, dedicate parts of their living space to flower gardens, pick wildflowers, or buy commercially-grown flowers from florists. Flower production and trade supports developing economies through their availability as a fair trade product.
Flowers provide less food than other major plant parts (seeds, fruits, roots, stems and leaves), but still provide several important vegetables and spices. Flower vegetables include broccoli, cauliflower and artichoke. The most expensive spice, saffron, consists of dried stigmas of a crocus. Other flower spices are cloves and capers. Hops flowers are used to flavor beer. Marigold flowers are fed to chickens to give their egg yolks a golden yellow color, which consumers find more desirable; dried and ground marigold flowers are also used as a spice and coloring agent in Georgian cuisine. Flowers of the dandelion and elder are often made into wine. Bee pollen, pollen collected from bees, is considered a health food by some people. Honey consists of bee-processed flower nectar and is often named for the type of flower, e.g. orange blossom honey, clover honey and tupelo honey.
Hundreds of fresh flowers are edible, but only few are widely marketed as food. They are often added to salads as garnishes. Squash blossoms are dipped in breadcrumbs and fried. Some edible flowers include nasturtium, chrysanthemum, carnation, cattail, Japanese honeysuckle, chicory, cornflower, canna, and sunflower. Edible flowers such as daisy, rose, and violet are sometimes candied.
Flowers such as chrysanthemum, rose, jasmine, Japanese honeysuckle, and chamomile, chosen for their fragrance and medicinal properties, are used as tisanes, either mixed with tea or on their own.
Flowers have been used since prehistoric times in funeral rituals: traces of pollen have been found on a woman's tomb in the El Miron Cave in Spain. Many cultures draw a connection between flowers and life and death, and because of their seasonal return flowers also suggest rebirth, which may explain why many people place flowers upon graves. The ancient Greeks, as recorded in Euripides's play The Phoenician Women, placed a crown of flowers on the head of the deceased; they also covered tombs with wreaths and flower petals. Flowers were widely used in ancient Egyptian burials, and the Mexicans to this day use flowers prominently in their Day of the Dead celebrations in the same way that their Aztec ancestors did.
Giving
The flower-giving tradition goes back to prehistoric times when flowers often had a medicinal and herbal attributes. Archaeologists found in several grave sites remnants of flower petals. Flowers were first used as sacrificial and burial objects. Ancient Egyptians and later Greeks and Romans used flowers. In Egypt, burial objects from the time around 1540 BC were found, which depicted red poppy, yellow Araun, cornflower and lilies. Records of flower giving appear in Chinese writings and Egyptian hieroglyphics, as well as in Greek and Roman mythology. The practice of giving a flower flourished in the Middle Ages when couples showed affection through flowers.
The tradition of flower-giving exists in many forms. It is an important part of Russian culture and folklore. It is common for students to give flowers to their teachers. To give yellow flowers in a romantic relationship means breakup in Russia. Nowadays, flowers are often given away in the form of a flower bouquet.
| Biology and health sciences | Biology | null |
4576951 | https://en.wikipedia.org/wiki/Convergence%20tests | Convergence tests | In mathematics, convergence tests are methods of testing for the convergence, conditional convergence, absolute convergence, interval of convergence or divergence of an infinite series .
List of tests
Limit of the summand
If the limit of the summand is undefined or nonzero, that is , then the series must diverge. In this sense, the partial sums are Cauchy only if this limit exists and is equal to zero. The test is inconclusive if the limit of the summand is zero. This is also known as the nth-term test, test for divergence, or the divergence test.
Ratio test
This is also known as d'Alembert's criterion.
Consider two limits and . If , the series diverges. If then the series converges absolutely. If then the test is inconclusive, and the series may converge absolutely, conditionally or diverge.
Root test
This is also known as the nth root test or Cauchy's criterion.
Let
where denotes the limit superior (possibly ; if the limit exists it is the same value).
If r < 1, then the series converges absolutely. If r > 1, then the series diverges. If r = 1, the root test is inconclusive, and the series may converge or diverge.
The root test is stronger than the ratio test: whenever the ratio test determines the convergence or divergence of an infinite series, the root test does too, but not conversely.
Integral test
The series can be compared to an integral to establish convergence or divergence. Let be a non-negative and monotonically decreasing function such that . If
then the series converges. But if the integral diverges, then the series does so as well.
In other words, the series converges if and only if the integral converges.
-series test
A commonly-used corollary of the integral test is the p-series test. Let . Then converges if .
The case of yields the harmonic series, which diverges. The case of is the Basel problem and the series converges to . In general, for , the series is equal to the Riemann zeta function applied to , that is .
Direct comparison test
If the series is an absolutely convergent series and for sufficiently large n , then the series converges absolutely.
Limit comparison test
If , (that is, each element of the two sequences is positive) and the limit exists, is finite and non-zero, then either both series converge or both series diverge.
Cauchy condensation test
Let be a non-negative non-increasing sequence. Then the sum converges if and only if the sum converges. Moreover, if they converge, then holds.
Abel's test
Suppose the following statements are true:
is a convergent series,
is a monotonic sequence, and
is bounded.
Then is also convergent.
Absolute convergence test
Every absolutely convergent series converges.
Alternating series test
Suppose the following statements are true:
is monotonic,
Then and are convergent series.
This test is also known as the Leibniz criterion.
Dirichlet's test
If is a sequence of real numbers and a sequence of complex numbers satisfying
for every positive integer N
where M is some constant, then the series
converges.
Cauchy's convergence test
A series is convergent if and only if for every there is a natural number N such that
holds for all and all .
Stolz–Cesàro theorem
Let and be two sequences of real numbers. Assume that is a strictly monotone and divergent sequence and the following limit exists:
Then, the limit
Weierstrass M-test
Suppose that (fn) is a sequence of real- or complex-valued functions defined on a set A, and that there is a sequence of non-negative numbers (Mn) satisfying the conditions
for all and all , and
converges.
Then the series
converges absolutely and uniformly on A.
Extensions to the ratio test
The ratio test may be inconclusive when the limit of the ratio is 1. Extensions to the ratio test, however, sometimes allows one to deal with this case.
Raabe–Duhamel's test
Let { an } be a sequence of positive numbers.
Define
If
exists there are three possibilities:
if L > 1 the series converges (this includes the case L = ∞)
if L < 1 the series diverges
and if L = 1 the test is inconclusive.
An alternative formulation of this test is as follows. Let } be a series of real numbers. Then if b > 1 and K (a natural number) exist such that
for all n > K then the series {an} is convergent.
Bertrand's test
Let { an } be a sequence of positive numbers.
Define
If
exists, there are three possibilities:
if L > 1 the series converges (this includes the case L = ∞)
if L < 1 the series diverges
and if L = 1 the test is inconclusive.
Gauss's test
Let { an } be a sequence of positive numbers. If for some β > 1, then converges if and diverges if .
Kummer's test
Let { an } be a sequence of positive numbers. Then:
(1) converges if and only if there is a sequence of positive numbers and a real number c > 0 such that .
(2) diverges if and only if there is a sequence of positive numbers such that
and diverges.
Abu-Mostafa's test
Let be an infinite series with real terms and let be any real function such that for all positive integers n and the second derivative exists at . Then converges absolutely if and diverges otherwise.
| Mathematics | Sequences and series | null |
4577392 | https://en.wikipedia.org/wiki/Edge%20cover | Edge cover | In graph theory, an edge cover of a graph is a set of edges such that every vertex of the graph is incident to at least one edge of the set.
In computer science, the minimum edge cover problem is the problem of finding an edge cover of minimum size. It is an optimization problem that belongs to the class of covering problems and can be solved in polynomial time.
Definition
Formally, an edge cover of a graph is a set of edges such that each vertex in is incident with at least one edge in . The set is said to cover the vertices of . The following figure shows examples of edge coverings in two graphs (the set is marked with red).
A minimum edge covering is an edge covering of smallest possible size. The edge covering number is the size of a minimum edge covering. The following figure shows examples of minimum edge coverings (again, the set is marked with red).
Note that the figure on the right is not only an edge cover but also a matching. In particular, it is a perfect matching: a matching in which every vertex is incident with exactly one edge in . A perfect matching (if it exists) is always a minimum edge covering.
Examples
The set of all edges is an edge cover, assuming that there are no degree-0 vertices.
The complete bipartite graph has edge covering number .
Algorithms
A smallest edge cover can be found in polynomial time by finding a maximum matching and extending it greedily so that all vertices are covered. In the following figure, a maximum matching is marked with red; the extra edges that were added to cover unmatched nodes are marked with blue. (The figure on the right shows a graph in which a maximum matching is a perfect matching; hence it already covers all vertices and no extra edges were needed.)
On the other hand, the related problem of finding a smallest vertex cover is an NP-hard problem.
Looking at the image it already becomes obvious why, for a given minimum edge cover and maximum matching , letting and be the number of edges in and respectively, we have: . Indeed, contains a maximum matching, so the edges of can be decomposed between the edges of a maximum matching, covering vertices, and the other edges that each cover one other vertex. Thus, as covers all of the vertices, we have giving the desired equality.
| Mathematics | Graph theory | null |
4577602 | https://en.wikipedia.org/wiki/Evolution%20of%20birds | Evolution of birds | The evolution of birds began in the Jurassic Period, with the earliest birds derived from a clade of theropod dinosaurs named Paraves. Birds are categorized as a biological class, Aves. For more than a century, the small theropod dinosaur Archaeopteryx lithographica from the Late Jurassic period was considered to have been the earliest bird. Modern phylogenies place birds in the dinosaur clade Theropoda. According to the current consensus, Aves and a sister group, the order Crocodilia, together are the sole living members of an unranked reptile clade, the Archosauria. Four distinct lineages of bird survived the Cretaceous–Paleogene extinction event 66 million years ago, giving rise to ostriches and relatives (Palaeognathae), waterfowl (Anseriformes), ground-living fowl (Galliformes), and "modern birds" (Neoaves).
Phylogenetically, Aves is usually defined as all descendants of the most recent common ancestor of a specific modern bird species (such as the house sparrow, Passer domesticus), and either Archaeopteryx, or some prehistoric species closer to Neornithes (to avoid the problems caused by the unclear relationships of Archaeopteryx to other theropods). If the latter classification is used then the larger group is termed Avialae. Currently, the relationship between non-avian dinosaurs, Archaeopteryx, and modern birds is still under debate.
Origins
There is significant evidence that birds emerged within theropod dinosaurs, specifically, that birds are members of Maniraptora, a group of theropods which includes dromaeosaurs and oviraptorids, among others. As more non-avian theropods that are closely related to birds are discovered, the formerly clear distinction between non-birds and birds becomes less so. This was noted in the 19th century, with Thomas Huxley writing:
We have had to stretch the definition of the class of birds so as to include birds with teeth and birds with paw-like fore limbs and long tails. There is no evidence that Compsognathus possessed feathers; but, if it did, it would be hard indeed to say whether it should be called a reptilian bird or an avian reptile.
Discoveries in northeast China (Liaoning Province) demonstrate that many small theropod dinosaurs did indeed have feathers, among them the compsognathid Sinosauropteryx and the microraptorian dromaeosaurid Sinornithosaurus. This has contributed to this ambiguity of where to draw the line between birds and reptiles. Cryptovolans, a dromaeosaurid found in 2002 (which may be a junior synonym of Microraptor) was capable of powered flight, possessing a sternal keel and ribs with uncinate processes. Cryptovolans seems to make a better "bird" than Archaeopteryx which lacks some of these modern bird features. Because some basal members of Dromaeosauridae, including Microraptor, were capable of powered flight, some paleontologists have suggested that dromaeosaurids are actually derived from a flying ancestor, and that the larger members became secondarily flightless, mirroring the loss of flight in modern paleognaths like the ostrich. The discoveries of further basal dromaeosaurids potentially capable of powered flight, such as Xiaotingia, has provided more evidence for the theory that flight was first developed in the bird line by early dromaeosaurids rather than later by Aves as was previously supposed.
Although ornithischian (bird-hipped) dinosaurs share the same hip structure as birds, birds actually originated from the saurischian (lizard-hipped) dinosaurs if the dinosaurian origin theory is correct. They thus arrived at their hip structure condition independently. In fact, a bird-like hip structure also developed a third time among a peculiar group of theropods, the Therizinosauridae.
An alternate theory to the dinosaurian origin of birds, espoused by a few scientists, notably Larry Martin and Alan Feduccia, states that birds (including maniraptoran "dinosaurs") evolved from early archosaurs like Longisquama. This theory is contested by most other paleontologists and experts in feather development and evolution.
Mesozoic birds
The basal bird Archaeopteryx, from the Jurassic, is well known as one of the first "missing links" to be found in support of evolution in the late 19th century. Though it is not considered a direct ancestor of modern birds, it gives a fair representation of how flight evolved and how the very first bird might have looked. It may be predated by Protoavis texensis, though the fragmentary nature of this fossil leaves it open to considerable doubt whether this was a bird ancestor. The skeleton of all early bird candidates is basically that of a small theropod dinosaur with long, clawed hands, though the exquisite preservation of the Solnhofen Plattenkalk shows Archaeopteryx was covered in feathers and had wings. While Archaeopteryx and its relatives may not have been very good fliers, they would at least have been competent gliders, setting the stage for the evolution of life on the wing.
The evolutionary trend among birds has been the reduction of anatomical elements to save weight. The first element to disappear was the bony tail, being reduced to a pygostyle and the tail function taken over by feathers. Confuciusornis is an example of their trend. While keeping the clawed fingers, perhaps for climbing, it had a pygostyle tail, though longer than in modern birds. A large group of birds, the Enantiornithes, evolved into ecological niches similar to those of modern birds and flourished throughout the Mesozoic. Though their wings resembled those of many modern bird groups, they retained the clawed wings and a snout with teeth rather than a beak in most forms. The loss of a long tail was followed by a rapid evolution of their legs which evolved to become highly versatile and adaptable tools that opened up new ecological niches.
The Cretaceous saw the rise of more modern birds with a more rigid ribcage with a carina and shoulders able to allow for a powerful upstroke, essential to sustained powered flight. Another improvement was the appearance of an alula, used to achieve better control of landing or flight at low speeds. They also had a more derived pygostyle, with a ploughshare-shaped end. An early example is Yanornis. Many were coastal birds, strikingly resembling modern shorebirds, like Ichthyornis, or ducks, like Gansus. Some evolved as swimming hunters, like the Hesperornithiformes – a group of flightless divers resembling grebes and loons. While modern in most respects, most of these birds retained typical reptilian-like teeth and sharp claws on the manus.
The modern toothless birds evolved from the toothed ancestors in the Cretaceous. Meanwhile, the earlier primitive birds, particularly the Enantiornithes, continued to thrive and diversify alongside the pterosaurs through this geologic period until they became extinct due to the K–T extinction event. All but a few groups of the toothless Neornithes were also cut short. The surviving lineages of birds were the comparatively primitive Palaeognathae (ostrich and its allies), the aquatic duck lineage, the terrestrial fowl, and the highly volant Neoaves.
Radiation of modern birds
Modern birds originated in the late Cretaceous. They are split into the paleognaths and neognaths. The paleognaths include the tinamous (grouse-like birds, found only in Central and South America) and the ratites, which nowadays are found almost exclusively in the Southern Hemisphere. The ratites are large flightless birds, and include ostriches, rheas, cassowaries, kiwis and emus. The ratites are a paraphyletic (artificial) grouping because tinamous are part of their evolutionary clade and they have likely lost the ability to fly independently, becoming an example of convergent evolution. However, the evidence about their evolution is still ambiguous, partly because there are no uncontroversial fossils from the Mesozoic and partly because their phylogenetic relationships are still uncertain.
The basal divergence within Neognathes is between Galloanserae and Neoaves.
The timing of divergence of these major groups are a matter of debate. It is agreed that modern birds originated in the Cretaceous and that the split between Galloanserae and Neoaves occurred before the Cretaceous–Paleogene extinction event, but there are different opinions about whether the radiation of the remaining neognaths occurred before or after the extinction event. This disagreement is in part caused by a divergence in the evidence, with molecular dating suggesting a Cretaceous radiation and the fossil record suggesting a Paleogene radiation. The latest attempts to reconcile the molecular and fossil evidence estimated the most recent common ancestor of modern birds at 95 million years ago and the split between Galloanseres and Neoaves at 85 million years ago. Notably, these studies show that the rapid proliferation of lineages in Neoaves seems to coincide with the Cretaceous–Paleogene extinction event, suggesting a role for ecological opportunity stimulating diversification in the aftermath of the mass extinction.
In contrast, another recent genomic study suggests that the Galloanserae and Neoaves diverged around the Early-Late Cretaceous boundary (100.5 million years ago), with the paleognaths and neognaths diverging even earlier (around 130 million years ago), and that most terrestrial neoavian orders gradually diverged from one another throughout the Late Cretaceous, roughly in sync with the concurrent radiation of flowering plants. This would suggest that a majority of all terrestrial avian orders coexisted with the non-avian dinosaurs and are K-Pg extinction survivors. In contrast, most major radiations of seabirds and shorebirds (as well as in paleognaths, despite their ancient origins) were found to have only occurred after the K-Pg extinction event, and primarily after the Paleocene–Eocene Thermal Maximum. This clashes with previous studies that found a very rapid radiation of avian orders only after the K-Pg extinction. The results of this study have been disputed by other researchers, due to a lack of fossil evidence to support its conclusions.
The birds that survived the end-of-Cretaceous extinction were likely ground-dwelling (not arboreal) and thus persisted despite the worldwide destruction of forests.
An analysis of the variation of diversification rates through time further revealed a potential effect of climate on the evolution diversification rates in birds in which the generation of new lineages accelerates during periods of global cooling. This can be the result of climate cooling fragmenting tropical biomes and producing widespread allopatric speciation plus an effect of some lineages diversifying in the expanding arid and cool biomes.
Bird skull evolution decelerated compared with the evolution of their dinosaur predecessors after the Cretaceous–Paleogene extinction event, rather than accelerating as often believed to have caused the cranial shape diversity of modern birds.
Classification of modern species
The phylogenetic classification of birds is a contentious issue. Sibley & Ahlquist's Phylogeny and Classification of Birds (1990) is a landmark work on the classification of birds (although frequently debated and constantly revised). A preponderance of evidence suggests that most modern bird orders constitute good clades. However, scientists are not in agreement as to the precise relationships between the main clades. Evidence from modern bird anatomy, fossils and DNA have all been brought to bear on the problem but no strong consensus has emerged.
Structural characteristics and fossil records have historically provided enough data for systematists to form hypotheses regarding the phylogenetic relationships between birds. Imprecisions within these methods is the main factor for why a lack of exact knowledge with regards to the orders and families of birds exists. Expansions in the study of computer-generated DNA sequencing and computer generated phylogenetics has provided a more accurate method for classifying bird species - although DNA data studying can only go so far, and questions are still unanswered.
Current evolutionary trends in birds
Evolution generally occurs at a scale far too slow to be witnessed by humans. However, bird species are currently going extinct at a far greater rate than any possible speciation or other generation of new species. The disappearance of a population, subspecies, or species represents the permanent loss of a range of genes.
Another concern with evolutionary implications is a suspected increase in hybridization. This may arise from human alteration of habitats enabling related allopatric species to overlap. Forest fragmentation can create extensive open areas, connecting previously isolated patches of open habitat. Populations that were isolated for sufficient time to diverge significantly, but not sufficient to be incapable of producing fertile offspring may now be interbreeding so broadly that the integrity of the original species may be compromised. For example, the many hybrid hummingbirds found in northwest South America may represent a threat to the conservation of the distinct species involved.
Several species of birds have been bred in captivity to create variations on wild species. In some birds this is limited to color variations, while others are bred for larger egg or meat production, for flightlessness or other characteristics.
In December 2019 the results of a joint study by Chicago's Field Museum and the University of Michigan into changes in the morphology of birds were published in Ecology Letters. The study uses bodies of birds which died as a result of colliding with buildings in Chicago, Illinois, since 1978. The sample is made up of over 70,000 specimens from 52 species and spans the period from 1978 to 2016. The study shows that the length of birds' lower leg bones (an indicator of body sizes) shortened by an average of 2.4% and their wings lengthened by 1.3%. The findings of the study suggest the morphological changes are the result of climate change, demonstrating an example of evolutionary change following Bergmann's rule.
| Biology and health sciences | Basics_4 | Biology |
4579350 | https://en.wikipedia.org/wiki/Continental%20margin | Continental margin | A continental margin is the outer edge of continental crust abutting oceanic crust under coastal waters. It is one of the three major zones of the ocean floor, the other two being deep-ocean basins and mid-ocean ridges. The continental margin consists of three different features: the continental rise, the continental slope, and the continental shelf. Continental margins constitute about 28% of the oceanic area.
Subzones
The continental shelf is the relatively shallow water area found in proximity to continents; it is the portion of the continental margin that transitions from the shore out towards to ocean. Continental shelves are believed to make up 7% of the sea floor. The width of continental shelves worldwide varies in the range of 0.03–1500 km. The continental shelf is generally flat, and ends at the shelf break, where there is a drastic increase in slope angle: The mean angle of continental shelves worldwide is 0° 07′, and typically steeper closer to the coastline than it is near the shelf break. At the shelf break begins the continental slope, which can be 1–5 km above the deep-ocean floor. The continental slope often exhibits features called submarine canyons. Submarine canyons often cut into the continental shelves deeply, with near vertical sides, and continue to cut the morphology to the abyssal plain. These canyons are often V-shaped, and can sometime enlarge onto the continental shelf. At the base of the continental slope, there is a sudden decrease in slope angle, and the sea floor begins to level out towards the abyssal plain. This portion of the seafloor is called the continental rise, and marks the outermost zone of the continental margin.
Types
There are two types of continental margins: active and passive margins.
Active margins are typically associated with lithospheric plate boundaries. These active margins can be convergent or transform margins, and are also places of high tectonic activity, including volcanoes and earthquakes. The West Coast of North America and South America are active margins. Active continental margins are typically narrow from coast to shelf break, with steep descents into trenches. Convergent active margins occur where oceanic plates meet continental plates. The denser oceanic crust of one plate subducts below the less dense continental crust of another plate. Convergent active margins are the most common type of active margin. Transform active margins are more rare, and occur when an oceanic plate and a continental plate are moving parallel to each other in opposite directions. These transform margins are often characterized by many offshore faults, which causes high degree of relief offshore, marked by islands, shallow banks, and deep basins. This is known as the continental borderland.
Passive margins are often located in the interior of lithospheric plates, away from the plate boundaries, and lack major tectonic activity. They often face mid-ocean ridges. From this, comes a wide variety of features, such as low-relief land extending miles away from the beach, long river systems and piles of sediment accumulating on the continental shelf. The East Coast of the United States is an example of a passive margin. These margins are much wider and less steep than active margins.
Sediment accumulation
As continental crust weathers and erodes, it degrades into mainly sands and clays. Many of these particles end up in streams and rivers that then dump into the ocean. Of all the sediment in the stream load, 80% is then trapped and dispersed on continental margins. While modern river sediment is often still preserved closer to shore, continental shelves show high levels of glacial and relict sediments, deposited when sea level was lower. Often found on passive margins are several kilometres of sediment, consisting of terrigenous and carbonate (biogenous) deposits. These sediment reservoirs are often useful in the study of paleoceanography and the original formation of ocean basins. These deposits are often not well preserved on active margin shelves due to tectonic activity.
Economic significance
The continental shelf is the most economically valuable part of the ocean. It often is the most productive portion of the continental margin, as well as the most studied portion, due to its relatively shallow, accessible depths.
Due to the rise of offshore drilling, mining and the limitations of fisheries off the continental shelf, the United Nations Convention on the Law of the Sea (UNCLOS) was established. The edge of the continental margin is one criterion for the boundary of the internationally recognized claims to underwater resources by countries in the definition of the "continental shelf" by the UNCLOS (although in the UN definition the "legal continental shelf" may extend beyond the geomorphological continental shelf and vice versa). Such resources include fishing grounds, oil and gas accumulations, sand, gravel, and some heavy minerals in the shallower areas of the margin. Metallic minerals resources are thought to also be associated with certain active margins, and of great value.
| Physical sciences | Tectonics | Earth science |
426604 | https://en.wikipedia.org/wiki/Vermilion | Vermilion | Vermilion (sometimes vermillion) is a color family and pigment most often used between antiquity and the 19th century from the powdered mineral cinnabar (a form of mercury sulfide). It is synonymous with red orange, which often takes a modern form, but is 11% brighter (at full brightness).
Etymology and common name
Used first in English in the 13th century, the word vermilion came from the Old French word vermeillon, which was derived from vermeil, from the Latin vermiculus the diminutive of the Latin word vermis for worm.
The name originated because it had a similar color to the natural red dye made from an insect, Kermes vermilio, which was widely used in Europe. The first recorded use of "vermilion" as a color name in English was in 1289.
The term cinnabar is used in mineralogy and crystallography for the red crystalline form of mercury sulfide HgS. Thus, the natural mineral pigment is called "cinnabar", and its synthetic form is called "vermilion" from red lead.
Chemistry and manufacture
Vermilion is a dense, opaque pigment with a clear, brilliant hue. The pigment was originally made by grinding a powder of cinnabar (mercury sulfide). Like most mercury compounds, it is toxic.
Vermilion is not one specific hue; mercuric sulfides make a range of warm hues, from bright orange-red to a duller reddish-purple that resembles fresh liver. Differences in hue are caused by the size of the ground particles of pigment. Larger crystals produce duller and less orange hues.
Cinnabar pigment was a side product of the mining of mercury, and mining cinnabar was difficult, expensive, and dangerous, because of the toxicity of mercury. Greek philosopher Theophrastus of Eresus (371–286 BC) described the process in De Lapidibus, the first scientific book on minerals. Efforts began early to find a better way to make the pigment.
The Chinese were probably the first to make a synthetic vermilion as early as the fourth century BC. Greek alchemist Zosimus of Panopolis (third–fourth century AD) wrote that such a method existed. In the early 9th century, the process was accurately described by Persian alchemist Jabir ibn Hayyan (722–804) in his book of recipes of colors, and the process began to be widely used in Europe.
The process described by Jabir ibn Hayyan was fairly simple:
Mix mercury with sulfur to form aethiopes mineralis, a black compound of mercury sulfide.
Heat this in a flask (the compound vaporizes and recondenses in the top of the flask).
Break the flask.
Collect the vermilion and grind it.
When first created, the material is almost black. As it is ground, the red color appears. The longer the compound is ground, the finer the color becomes. Italian Renaissance artist Cennino Cennini wrote: "If you were to grind it every day, even for 20 years, it would keep getting better and more perfect."
In the 17th century, a new method of making the pigment was introduced, known as the Dutch method. Mercury and melted sulfur were mashed to make black mercury sulfide, then heated in a retort, producing vapors condensing as a bright, red mercury sulfide. To remove the sulfur, these crystals were treated with a strong alkali, washed, and finally ground under water to yield the commercial powder form of the pigment. The pigment is still made today using essentially the same process.
Vermilion has one important defect; it is liable to darken, or develop a purplish-gray surface sheen. Cennino Cennini wrote, "Bear in mind ... that it is not in its character to be exposed to air, but it is more resistant on panel than on walls since, when it is used and laid on a wall, over a period of time, standing in the air, it turns black." Newer research indicates that chlorine ions and light may aid in decomposing vermilion into elemental mercury, which is black when in finely dispersed form.
Vermilion was the primary red pigment used by European painters, from the Renaissance until the 20th century. Because of its cost and toxicity, though, it was almost entirely replaced by a new synthetic pigment, cadmium red, in the 20th century. As cadmium can also be toxic, some scientists propose replacing this with solid solutions of the perovskites CaTaO2N and LaTaON2.
Genuine vermilion pigment today comes mostly from China; it is a synthetic mercuric sulfide, labeled on paint tubes as PR-106 (Red Pigment 106). The synthetic pigment is of higher quality than vermilion made from ground cinnabar, which has many impurities. The pigment is very toxic, and should be used with great care.
Gallery
History
The colors are widely used in the art and decoration of Ancient Rome and the Byzantine Empire, then in the illuminated manuscripts of the Middle Ages, in the paintings of the Renaissance, and in the art and lacquerware of China.
Antiquity
The first documented use of vermilion pigment, made with ground cinnabar, dates to 8000–7000 BC, and was found at the neolithic village of Çatalhöyük, in modern-day Turkey. Cinnabar was mined in Spain beginning in about 5300 BC. In China, the first documented use of cinnabar as a pigment was by the Yangshao culture (5000–4000 BC), where it was used to paint ceramics, to cover the walls and floors of rooms, and for ritual ceremonies.
The principal source of cinnabar for the ancient Romans was the Almaden mine in northwest Spain, which was worked by prisoners. Since the ore of mercury was highly toxic, a term in the mines was a near-guaranteed death sentence. Pliny the Elder described the mines this way:
Nothing is more carefully guarded. It is forbidden to break up or refine the cinnabar on the spot. They send it to Rome in its natural condition, under seal, to the extent of some ten thousand librae (Roman pounds thus 3289 kg) a year. The sales price is fixed by law to keep it from becoming impossibly expensive, and the price fixed is seventy sesterces a pound.
In Rome, the precious pigment was used to paint frescoes, decorate statues, and even as a cosmetic. In Roman triumphs, the victors had their faces covered with vermilion powder, and the face of Jupiter on the Capitoline Hill was also colored vermilion. Cinnabar was used to paint the walls of some of the most luxurious villas in Pompeii, including the Villa of the Mysteries (Italian: Villa dei Misteri). Pliny reported its painters stole a large portion of the expensive pigment by frequently washing their brushes and saving the wash water.
In the Byzantine Empire, the use of cinnabar/the vermilion color was reserved for the use of the imperial family and administrators; official letters and imperial decrees were written in vermilion ink, made with cinnabar.
In South Asia
It is known as sindoor. Sindoor is commonly used by married women in Hindu religion.
In the Americas
Vermilion was also used by the native peoples of America, to paint ceramics, figurines, and murals, and for the decoration of burials. It was used in the Chavin civilization (400 BC – 200 AD), and in the Maya, Sican, Moche, and Inca empires. The major source was the Huancavelica mine in the Andes mountains in central Peru.
The most dramatic example of vermilion use in the Americas was the so-called Tomb of the Red Queen, located in Temple XIII, in the ruins of the Mayan city of Palenque in Chiapas, Mexico. The temple is dated to between 600 and 700 AD. It was discovered in 1994 by Mexican archeologist Fanny López Jiménez. The body and all objects in the sarcophagus were covered with bright red vermilion powder made from cinnabar.
In the Middle Ages and Renaissance
The technique for making a synthetic vermilion by combining sulfur and mercury was in use in Europe in the 9th century, but the pigment was still expensive. Since it was almost as expensive as gold leaf, it was used only in the most important decoration of illuminated manuscripts, while the less expensive minium, made with red lead, was used for the red letters and symbols in the text.
Vermilion was also used by painters in the Renaissance as a very vivid and bright red, though it did have the weakness of sometimes turning dark with time. Florentine artist Cennino Cennini described it in his handbook for artists:
By the 20th century, the cost and toxicity of vermilion led to its gradually being replaced by synthetic pigments, particularly cadmium red, which had a comparable color and opacity.
Chinese red
In China, the color vermilion was also playing an important role in national culture. The color was mostly used in creating Chinese lacquerware, which was exported around the world, giving rise to the term "Chinese red".
The lacquer came from the Chinese lacquer tree, or Toxicodendron vernicifluum, a relative of the sumac tree, which grew in regions of China, Korea, and Japan. The sap or resin of the tree, called urushiol, was caustic and toxic (it contained the same chemical compound as poison ivy), but painted onto wood or metal, it hardened into a fine natural plastic, or lacquer surface. The pure sap was dark brown, but beginning in about the third century BC, during the Han dynasty, Chinese artisans colored it with powdered cinnabar or with red ochre (ferric oxide), giving it an orange-red color. Beginning in about the 8th century, Chinese chemists began making synthetic vermilion from mercury and sulfur, which reduced the price of the pigment and allowed the production of Chinese lacquerware on a larger scale.
The shade of red of the lacquerware has changed over the centuries. During the Eastern Han dynasty (25–220 AD) the Chinese word for red referred to a light red. However, during the Tang dynasty (618–907), when the synthetic vermilion was introduced, that color became darker and richer. The poet Bai Juyi (772–846) wrote in a song poem praising Jiangnan, "the flowers by the river when the sun rises are redder than flames", and the word he used for red was the word for vermilion, or Chinese red.
When Chinese lacquerware and the ground cinnabar used to color it were exported to Europe in the 17th and 18th centuries, European collectors considered it to be finer than the European vermilion. In 1835, "Chinese vermilion" was described as a cinnabar so pure that it only had to be ground into powder to become a perfect vermilion. Historically, European vermilion often included adulterants including brick, orpiment, iron oxide, Persian red, iodine scarlet—and minium (red lead), an inexpensive and bright, but fugitive lead-oxide pigment.
Since ancient times, vermilion was regarded as the color of blood, thus the color of life. It was used to paint temples and the carriages of the emperor, and as the printing paste for personal seals. It was also used for unique red calligraphic ink reserved for emperors. Chinese Taoists associated vermilion with eternity.
In nature
Vermilion flycatcher
Vermilion cardinal
Vermilion tanager
In art and culture
Religion
The Shaolin temple, where Buddhist monk Bodhidharma is reputed to have established the new sect of Chan Buddhism (Zen Buddhism), is colored a bright tone of vermilion. This temple was featured in the West by the 1972–1975 TV series Kung Fu.
In the Bible, vermilion is listed as a pigment that was in use for painting buildings during the reign of Shallum the son of Josiah king of Judah, and is named in the book of the prophet Ezekiel as a pigment used in art that portrayed Chaldean men. (Jeremiah 22:11–14, Ezekiel 23:14–17)
The vermilion rose is a symbol of the Blessed Virgin Mary.
Hindu women use vermilion along the hair parting line known as sindoor, to signify that they are married. Hindu men and women often wear vermilion on their foreheads during religious ceremonies and festivals.
In Shintoism, Torii Gates which mark the entrances to sacred spaces, as well as the columns and eaves of shrines, are traditionally painted vermilion to ward off evil.
Mythology
In Han China's Five Elements cosmology (cf. Chinese mythology), one of the four symbols of the four directions is a bird called Vermilion Bird, which represents the direction of south. The color red (particularly as exemplified by cinnabar/vermilion) was also symbolically associated with summer, fire, a certain note on the musical scale, a certain day of the calendar, etc.
Literature
Vermilion Sands is a collection of science-fiction short stories by J. G. Ballard published in 1971 about an imaginary future resort that pleases its guests by using various kinds of futuristic technology.
Manfred, a short dramatic poem by Lord Byron: "...With the azure and vermillion / which is mixed for my pavilion"
Music
The Dutch singer Simone Simons, released her debut solo album called Vermillion in 2024. On this album, the word vermillion was used for one of the song titles called "Vermillion Dreams" and also the album was called vermillion.
Video games
Vermilion City is one of the locations used in the English-translated versions of the Pokémon video games and anime. It is a port town in the Kanto area, and the name is derived from the original Japanese name クチバシティ (Kuchiba City). Kuchiba is an orange-red color associated with sunsets and autumnal leaves and "Vermilion" was used as an approximate translation.
Variations
Red-orange
The Crayola color red-orange has been a Crayola color since 1930.
Orange-red
The web color orange-red was formulated in 1987 as one of the X11 colors, which became known as the X11 web colors after the invention of the World Wide Web in 1991.
Medium vermilion
This color is the medium tone of vermilion called vermilion on the Plochere color list, which was formulated in 1948 and is used widely by interior designers.
Chinese red
Chinese red or China red is the name used for the vermilion shade used in Chinese lacquerware. The shade of the color can vary from dark to light depending upon how the pigment is made and how the lacquer was applied. Chinese red was originally made from the powdered mineral cinnabar, but beginning in about the 8th century it was made more commonly by a chemical process combining mercury and sulfur. Vermilion has significance in Taoist culture, and is regarded as the color of life and eternity.
"Chinese red" appears in English in 1924.
| Physical sciences | Colors | Physics |
426743 | https://en.wikipedia.org/wiki/Graph%20coloring | Graph coloring | In graph theory, graph coloring is a methodic assignment of labels traditionally called "colors" to elements of a graph. The assignment is subject to certain constraints, such as that no two adjacent elements have the same color. Graph coloring is a special case of graph labeling. In its simplest form, it is a way of coloring the vertices of a graph such that no two adjacent vertices are of the same color; this is called a vertex coloring. Similarly, an edge coloring assigns a color to each edges so that no two adjacent edges are of the same color, and a face coloring of a planar graph assigns a color to each face (or region) so that no two faces that share a boundary have the same color.
Vertex coloring is often used to introduce graph coloring problems, since other coloring problems can be transformed into a vertex coloring instance. For example, an edge coloring of a graph is just a vertex coloring of its line graph, and a face coloring of a plane graph is just a vertex coloring of its dual. However, non-vertex coloring problems are often stated and studied as-is. This is partly pedagogical, and partly because some problems are best studied in their non-vertex form, as in the case of edge coloring.
The convention of using colors originates from coloring the countries in a political map, where each face is literally colored. This was generalized to coloring the faces of a graph embedded in the plane. By planar duality it became coloring the vertices, and in this form it generalizes to all graphs. In mathematical and computer representations, it is typical to use the first few positive or non-negative integers as the "colors". In general, one can use any finite set as the "color set". The nature of the coloring problem depends on the number of colors but not on what they are.
Graph coloring enjoys many practical applications as well as theoretical challenges. Beside the classical types of problems, different limitations can also be set on the graph, or on the way a color is assigned, or even on the color itself. It has even reached popularity with the general public in the form of the popular number puzzle Sudoku. Graph coloring is still a very active field of research.
History
The first results about graph coloring deal almost exclusively with planar graphs in the form of map coloring.
While trying to color a map of the counties of England, Francis Guthrie postulated the four color conjecture, noting that four colors were sufficient to color the map so that no regions sharing a common border received the same color. Guthrie's brother passed on the question to his mathematics teacher Augustus De Morgan at University College, who mentioned it in a letter to William Hamilton in 1852. Arthur Cayley raised the problem at a meeting of the London Mathematical Society in 1879. The same year, Alfred Kempe published a paper that claimed to establish the result, and for a decade the four color problem was considered solved. For his accomplishment Kempe was elected a Fellow of the Royal Society and later President of the London Mathematical Society.
In 1890, Percy John Heawood pointed out that Kempe's argument was wrong. However, in that paper he proved the five color theorem, saying that every planar map can be colored with no more than five colors, using ideas of Kempe. In the following century, a vast amount of work was done and theories were developed to reduce the number of colors to four, until the four color theorem was finally proved in 1976 by Kenneth Appel and Wolfgang Haken. The proof went back to the ideas of Heawood and Kempe and largely disregarded the intervening developments.
The proof of the four color theorem is noteworthy, aside from its solution of a century-old problem, for being the first major computer-aided proof.
In 1912, George David Birkhoff introduced the chromatic polynomial to study the coloring problem, which was generalised to the Tutte polynomial by W. T. Tutte, both of which are important invariants in algebraic graph theory. Kempe had already drawn attention to the general, non-planar case in 1879, and many results on generalisations of planar graph coloring to surfaces of higher order followed in the early 20th century.
In 1960, Claude Berge formulated another conjecture about graph coloring, the strong perfect graph conjecture, originally motivated by an information-theoretic concept called the zero-error capacity of a graph introduced by Shannon. The conjecture remained unresolved for 40 years, until it was established as the celebrated strong perfect graph theorem by Chudnovsky, Robertson, Seymour, and Thomas in 2002.
Graph coloring has been studied as an algorithmic problem since the early 1970s: the chromatic number problem (see section below) is one of Karp's 21 NP-complete problems from 1972, and at approximately the same time various exponential-time algorithms were developed based on backtracking and on the deletion-contraction recurrence of . One of the major applications of graph coloring, register allocation in compilers, was introduced in 1981.
Definition and terminology
Vertex coloring
When used without any qualification, a coloring of a graph almost always refers to a proper vertex coloring, namely a labeling of the graph's vertices with colors such that no two vertices sharing the same edge have the same color. Since a vertex with a loop (i.e. a connection directly back to itself) could never be properly colored, it is understood that graphs in this context are loopless.
The terminology of using colors for vertex labels goes back to map coloring. Labels like red and blue are only used when the number of colors is small, and normally it is understood that the labels are drawn from the integers .
A coloring using at most colors is called a (proper) -coloring. The smallest number of colors needed to color a graph is called its chromatic number, and is often denoted . Sometimes is used, since is also used to denote the Euler characteristic of a graph. A graph that can be assigned a (proper) -coloring is -colorable, and it is -chromatic if its chromatic number is exactly . A subset of vertices assigned to the same color is called a color class, every such class forms an independent set. Thus, a -coloring is the same as a partition of the vertex set into independent sets, and the terms -partite and -colorable have the same meaning.
Chromatic polynomial
The chromatic polynomial counts the number of ways a graph can be colored using some of a given number of colors. For example, using three colors, the graph in the adjacent image can be colored in 12 ways. With only two colors, it cannot be colored at all. With four colors, it can be colored in 24 + 4 × 12 = 72 ways: using all four colors, there are 4! = 24 valid colorings (every assignment of four colors to any 4-vertex graph is a proper coloring); and for every choice of three of the four colors, there are 12 valid 3-colorings. So, for the graph in the example, a table of the number of valid colorings would start like this:
The chromatic polynomial is a function that counts the number of -colorings of . As the name indicates, for a given the function is indeed a polynomial in . For the example graph, , and indeed .
The chromatic polynomial includes more information about the colorability of than does the chromatic number. Indeed, is the smallest positive integer that is not a zero of the chromatic polynomial .
Edge coloring
An edge coloring of a graph is a proper coloring of the edges, meaning an assignment of colors to edges so that no vertex is incident to two edges of the same color. An edge coloring with colors is called a -edge-coloring and is equivalent to the problem of partitioning the edge set into matchings. The smallest number of colors needed for an edge coloring of a graph is the chromatic index, or edge chromatic number, . A Tait coloring is a 3-edge coloring of a cubic graph. The four color theorem is equivalent to the assertion that every planar cubic bridgeless graph admits a Tait coloring.
Total coloring
Total coloring is a type of coloring on the vertices and edges of a graph. When used without any qualification, a total coloring is always assumed to be proper in the sense that no adjacent vertices, no adjacent edges, and no edge and its end-vertices are assigned the same color. The total chromatic number of a graph is the fewest colors needed in any total coloring of .
Face coloring
For a graph with a strong embedding on a surface, the face coloring is the dual of the vertex coloring problem.
Tutte's flow theory
For a graph G with a strong embedding on an orientable surface, William T. Tutte discovered that if the graph is k-face-colorable then G admits a nowhere-zero k-flow. The equivalence holds if the surface is sphere.
Unlabeled coloring
An unlabeled coloring of a graph is an orbit of a coloring under the action of the automorphism group of the graph. The colors remain labeled; it is the graph that is unlabeled.
There is an analogue of the chromatic polynomial which counts the number of unlabeled colorings of a graph from a given finite color set.
If we interpret a coloring of a graph on vertices as a vector in , the action of an automorphism is a permutation of the coefficients in the coloring vector.
Properties
Upper bounds on the chromatic number
Assigning distinct colors to distinct vertices always yields a proper coloring, so
The only graphs that can be 1-colored are edgeless graphs. A complete graph of n vertices requires colors. In an optimal coloring there must be at least one of the graph's m edges between every pair of color classes, so
More generally a family of graphs is -bounded if there is some function such that the graphs in can be colored with at most colors, where is the clique number of . For the family of the perfect graphs this function is .
The 2-colorable graphs are exactly the bipartite graphs, including trees and forests.
By the four color theorem, every planar graph can be 4-colored.
A greedy coloring shows that every graph can be colored with one more color than the maximum vertex degree,
Complete graphs have and , and odd cycles have and , so for these graphs this bound is best possible. In all other cases, the bound can be slightly improved; Brooks' theorem states that
Brooks' theorem: for a connected, simple graph G, unless G is a complete graph or an odd cycle.
Lower bounds on the chromatic number
Several lower bounds for the chromatic bounds have been discovered over the years:
If G contains a clique of size k, then at least k colors are needed to color that clique; in other words, the chromatic number is at least the clique number:
For perfect graphs this bound is tight. Finding cliques is known as the clique problem.
Hoffman's bound: Let be a real symmetric matrix such that whenever is not an edge in . Define , where are the largest and smallest eigenvalues of . Define , with as above. Then:
: Let be a positive semi-definite matrix such that whenever is an edge in . Define to be the least k for which such a matrix exists. Then
Lovász number: The Lovász number of a complementary graph is also a lower bound on the chromatic number:
Fractional chromatic number: The fractional chromatic number of a graph is a lower bound on the chromatic number as well:
These bounds are ordered as follows:
Graphs with high chromatic number
Graphs with large cliques have a high chromatic number, but the opposite is not true. The Grötzsch graph is an example of a 4-chromatic graph without a triangle, and the example can be generalized to the Mycielskians.
Theorem (, , ): There exist triangle-free graphs with arbitrarily high chromatic number.
To prove this, both, Mycielski and Zykov, each gave a construction of an inductively defined family of triangle-free graphs but with arbitrarily large chromatic number. constructed axis aligned boxes in whose intersection graph is triangle-free and requires arbitrarily many colors to be properly colored. This family of graphs is then called the Burling graphs. The same class of graphs is used for the construction of a family of triangle-free line segments in the plane, given by Pawlik et al. (2014). It shows that the chromatic number of its intersection graph is arbitrarily large as well. Hence, this implies that axis aligned boxes in as well as line segments in are not χ-bounded.
From Brooks's theorem, graphs with high chromatic number must have high maximum degree. But colorability is not an entirely local phenomenon: A graph with high girth looks locally like a tree, because all cycles are long, but its chromatic number need not be 2:
Theorem (Erdős): There exist graphs of arbitrarily high girth and chromatic number.
Bounds on the chromatic index
An edge coloring of G is a vertex coloring of its line graph , and vice versa. Thus,
There is a strong relationship between edge colorability and the graph's maximum degree . Since all edges incident to the same vertex need their own color, we have
Moreover,
Kőnig's theorem: if G is bipartite.
In general, the relationship is even stronger than what Brooks's theorem gives for vertex coloring:
Vizing's Theorem: A graph of maximal degree has edge-chromatic number or .
Other properties
A graph has a k-coloring if and only if it has an acyclic orientation for which the longest path has length at most k; this is the Gallai–Hasse–Roy–Vitaver theorem .
For planar graphs, vertex colorings are essentially dual to nowhere-zero flows.
About infinite graphs, much less is known.
The following are two of the few results about infinite graph coloring:
If all finite subgraphs of an infinite graph G are k-colorable, then so is G, under the assumption of the axiom of choice. This is the de Bruijn–Erdős theorem of .
If a graph admits a full n-coloring for every n ≥ n0, it admits an infinite full coloring .
Open problems
As stated above, A conjecture of Reed from 1998 is that the value is essentially closer to the lower bound,
The chromatic number of the plane, where two points are adjacent if they have unit distance, is unknown, although it is one of 5, 6, or 7. Other open problems concerning the chromatic number of graphs include the Hadwiger conjecture stating that every graph with chromatic number k has a complete graph on k vertices as a minor, the Erdős–Faber–Lovász conjecture bounding the chromatic number of unions of complete graphs that have at most one vertex in common to each pair, and the Albertson conjecture that among k-chromatic graphs the complete graphs are the ones with smallest crossing number.
When Birkhoff and Lewis introduced the chromatic polynomial in their attack on the four-color theorem, they conjectured that for planar graphs G, the polynomial has no zeros in the region . Although it is known that such a chromatic polynomial has no zeros in the region and that , their conjecture is still unresolved. It also remains an unsolved problem to characterize graphs which have the same chromatic polynomial and to determine which polynomials are chromatic.
Algorithms
Polynomial time
Determining if a graph can be colored with 2 colors is equivalent to determining whether or not the graph is bipartite, and thus computable in linear time using breadth-first search or depth-first search. More generally, the chromatic number and a corresponding coloring of perfect graphs can be computed in polynomial time using semidefinite programming. Closed formulas for chromatic polynomials are known for many classes of graphs, such as forests, chordal graphs, cycles, wheels, and ladders, so these can be evaluated in polynomial time.
If the graph is planar and has low branch-width (or is nonplanar but with a known branch decomposition), then it can be solved in polynomial time using dynamic programming. In general, the time required is polynomial in the graph size, but exponential in the branch-width.
Exact algorithms
Brute-force search for a k-coloring considers each of the assignments of k colors to n vertices and checks for each if it is legal. To compute the chromatic number and the chromatic polynomial, this procedure is used for every , impractical for all but the smallest input graphs.
Using dynamic programming and a bound on the number of maximal independent sets, k-colorability can be decided in time and space . Using the principle of inclusion–exclusion and Yates's algorithm for the fast zeta transform, k-colorability can be decided in time for any k. Faster algorithms are known for 3- and 4-colorability, which can be decided in time and , respectively. Exponentially faster algorithms are also known for 5- and 6-colorability, as well as for restricted families of graphs, including sparse graphs.
Contraction
The contraction of a graph G is the graph obtained by identifying the vertices u and v, and removing any edges between them. The remaining edges originally incident to u or v are now incident to their identification (i.e., the new fused node uv). This operation plays a major role in the analysis of graph coloring.
The chromatic number satisfies the recurrence relation:
due to , where u and v are non-adjacent vertices, and is the graph with the edge added. Several algorithms are based on evaluating this recurrence and the resulting computation tree is sometimes called a Zykov tree. The running time is based on a heuristic for choosing the vertices u and v.
The chromatic polynomial satisfies the following recurrence relation
where u and v are adjacent vertices, and is the graph with the edge removed. represents the number of possible proper colorings of the graph, where the vertices may have the same or different colors. Then the proper colorings arise from two different graphs. To explain, if the vertices u and v have different colors, then we might as well consider a graph where u and v are adjacent. If u and v have the same colors, we might as well consider a graph where u and v are contracted. Tutte's curiosity about which other graph properties satisfied this recurrence led him to discover a bivariate generalization of the chromatic polynomial, the Tutte polynomial.
These expressions give rise to a recursive procedure called the deletion–contraction algorithm, which forms the basis of many algorithms for graph coloring. The running time satisfies the same recurrence relation as the Fibonacci numbers, so in the worst case the algorithm runs in time within a polynomial factor of for n vertices and m edges. The analysis can be improved to within a polynomial factor of the number of spanning trees of the input graph. In practice, branch and bound strategies and graph isomorphism rejection are employed to avoid some recursive calls. The running time depends on the heuristic used to pick the vertex pair.
Greedy coloring
The greedy algorithm considers the vertices in a specific order , ..., and assigns to the smallest available color not used by 's neighbours among , ..., , adding a fresh color if needed. The quality of the resulting coloring depends on the chosen ordering. There exists an ordering that leads to a greedy coloring with the optimal number of colors. On the other hand, greedy colorings can be arbitrarily bad; for example, the crown graph on n vertices can be 2-colored, but has an ordering that leads to a greedy coloring with colors.
For chordal graphs, and for special cases of chordal graphs such as interval graphs and indifference graphs, the greedy coloring algorithm can be used to find optimal colorings in polynomial time, by choosing the vertex ordering to be the reverse of a perfect elimination ordering for the graph. The perfectly orderable graphs generalize this property, but it is NP-hard to find a perfect ordering of these graphs.
If the vertices are ordered according to their degrees, the resulting greedy coloring uses at most colors, at most one more than the graph's maximum degree. This heuristic is sometimes called the Welsh–Powell algorithm. Another heuristic due to Brélaz establishes the ordering dynamically while the algorithm proceeds, choosing next the vertex adjacent to the largest number of different colors. Many other graph coloring heuristics are similarly based on greedy coloring for a specific static or dynamic strategy of ordering the vertices, these algorithms are sometimes called sequential coloring algorithms.
The maximum (worst) number of colors that can be obtained by the greedy algorithm, by using a vertex ordering chosen to maximize this number, is called the Grundy number of a graph.
Heuristic algorithms
Two well-known polynomial-time heuristics for graph colouring are the DSatur and recursive largest first (RLF) algorithms.
Similarly to the greedy colouring algorithm, DSatur colours the vertices of a graph one after another, expending a previously unused colour when needed. Once a new vertex has been coloured, the algorithm determines which of the remaining uncoloured vertices has the highest number of different colours in its neighbourhood and colours this vertex next. This is defined as the degree of saturation of a given vertex.
The recursive largest first algorithm operates in a different fashion by constructing each color class one at a time. It does this by identifying a maximal independent set of vertices in the graph using specialised heuristic rules. It then assigns these vertices to the same color and removes them from the graph. These actions are repeated on the remaining subgraph until no vertices remain.
The worst-case complexity of DSatur is , where is the number of vertices in the graph. The algorithm can also be implemented using a binary heap to store saturation degrees, operating in where is the number of edges in the graph. This produces much faster runs with sparse graphs. The overall complexity of RLF is slightly higher than DSatur at .
DSatur and RLF are exact for bipartite, cycle, and wheel graphs.
Parallel and distributed algorithms
It is known that a -chromatic graph can be -colored in the deterministic LOCAL model, in . rounds, with . A matching lower bound of rounds is also known. This lower bound holds even if quantum computers that can exchange quantum information, possibly with a pre-shared entangled state, are allowed.
In the field of distributed algorithms, graph coloring is closely related to the problem of symmetry breaking. The current state-of-the-art randomized algorithms are faster for sufficiently large maximum degree Δ than deterministic algorithms. The fastest randomized algorithms employ the multi-trials technique by Schneider and Wattenhofer.
In a symmetric graph, a deterministic distributed algorithm cannot find a proper vertex coloring. Some auxiliary information is needed in order to break symmetry. A standard assumption is that initially each node has a unique identifier, for example, from the set . Put otherwise, we assume that we are given an n-coloring. The challenge is to reduce the number of colors from n to, e.g., Δ + 1. The more colors are employed, e.g. O(Δ) instead of Δ + 1, the fewer communication rounds are required.
A straightforward distributed version of the greedy algorithm for (Δ + 1)-coloring requires Θ(n) communication rounds in the worst case – information may need to be propagated from one side of the network to another side.
The simplest interesting case is an n-cycle. Richard Cole and Uzi Vishkin show that there is a distributed algorithm that reduces the number of colors from n to O(log n) in one synchronous communication step. By iterating the same procedure, it is possible to obtain a 3-coloring of an n-cycle in O( n) communication steps (assuming that we have unique node identifiers).
The function , iterated logarithm, is an extremely slowly growing function, "almost constant". Hence the result by Cole and Vishkin raised the question of whether there is a constant-time distributed algorithm for 3-coloring an n-cycle. showed that this is not possible: any deterministic distributed algorithm requires Ω( n) communication steps to reduce an n-coloring to a 3-coloring in an n-cycle.
The technique by Cole and Vishkin can be applied in arbitrary bounded-degree graphs as well; the running time is poly(Δ) + O( n). The technique was extended to unit disk graphs by Schneider and Wattenhofer. The fastest deterministic algorithms for (Δ + 1)-coloring for small Δ are due to Leonid Barenboim, Michael Elkin and Fabian Kuhn. The algorithm by Barenboim et al. runs in time O(Δ) + (n)/2, which is optimal in terms of n since the constant factor 1/2 cannot be improved due to Linial's lower bound. use network decompositions to compute a Δ+1 coloring in time .
The problem of edge coloring has also been studied in the distributed model. achieve a (2Δ − 1)-coloring in O(Δ + n) time in this model. The lower bound for distributed vertex coloring due to applies to the distributed edge coloring problem as well.
Decentralized algorithms
Decentralized algorithms are ones where no message passing is allowed (in contrast to distributed algorithms where local message passing takes places), and efficient decentralized algorithms exist that will color a graph if a proper coloring exists. These assume that a vertex is able to sense whether any of its neighbors are using the same color as the vertex i.e., whether a local conflict exists. This is a mild assumption in many applications e.g. in wireless channel allocation it is usually reasonable to assume that a station will be able to detect whether other interfering transmitters are using the same channel (e.g. by measuring the SINR). This sensing information is sufficient to allow algorithms based on learning automata to find a proper graph coloring with probability one.
Computational complexity
Graph coloring is computationally hard. It is NP-complete to decide if a given graph admits a k-coloring for a given k except for the cases k ∈ . In particular, it is NP-hard to compute the chromatic number. The 3-coloring problem remains NP-complete even on 4-regular planar graphs. On graphs with maximal degree 3 or less, however, Brooks' theorem implies that the 3-coloring problem can be solved in linear time. Further, for every k > 3, a k-coloring of a planar graph exists by the four color theorem, and it is possible to find such a coloring in polynomial time. However, finding the lexicographically smallest 4-coloring of a planar graph is NP-complete.
The best known approximation algorithm computes a coloring of size at most within a factor O(n(log log n)2(log n)−3) of the chromatic number. For all ε > 0, approximating the chromatic number within n1−ε is NP-hard.
It is also NP-hard to color a 3-colorable graph with 5 colors, 4-colorable graph with 7 colours, and a k-colorable graph with colors for k ≥ 5.
Computing the coefficients of the chromatic polynomial is ♯P-hard. In fact, even computing the value of is ♯P-hard at any rational point k except for k = 1 and k = 2. There is no FPRAS for evaluating the chromatic polynomial at any rational point k ≥ 1.5 except for k = 2 unless NP = RP.
For edge coloring, the proof of Vizing's result gives an algorithm that uses at most Δ+1 colors. However, deciding between the two candidate values for the edge chromatic number is NP-complete. In terms of approximation algorithms, Vizing's algorithm shows that the edge chromatic number can be approximated to within 4/3,
and the hardness result shows that no (4/3 − ε)-algorithm exists for any ε > 0 unless P = NP. These are among the oldest results in the literature of approximation algorithms, even though neither paper makes explicit use of that notion.
Applications
Scheduling
Vertex coloring models to a number of scheduling problems. In the cleanest form, a given set of jobs need to be assigned to time slots, each job requires one such slot. Jobs can be scheduled in any order, but pairs of jobs may be in conflict in the sense that they may not be assigned to the same time slot, for example because they both rely on a shared resource. The corresponding graph contains a vertex for every job and an edge for every conflicting pair of jobs. The chromatic number of the graph is exactly the minimum makespan, the optimal time to finish all jobs without conflicts.
Details of the scheduling problem define the structure of the graph. For example, when assigning aircraft to flights, the resulting conflict graph is an interval graph, so the coloring problem can be solved efficiently. In bandwidth allocation to radio stations, the resulting conflict graph is a unit disk graph, so the coloring problem is 3-approximable.
Register allocation
A compiler is a computer program that translates one computer language into another. To improve the execution time of the resulting code, one of the techniques of compiler optimization is register allocation, where the most frequently used values of the compiled program are kept in the fast processor registers. Ideally, values are assigned to registers so that they can all reside in the registers when they are used.
The textbook approach to this problem is to model it as a graph coloring problem. The compiler constructs an interference graph, where vertices are variables and an edge connects two vertices if they are needed at the same time. If the graph can be colored with k colors then any set of variables needed at the same time can be stored in at most k registers.
Other applications
The problem of coloring a graph arises in many practical areas such as sports scheduling, designing seating plans, exam timetabling, the scheduling of taxis, and solving Sudoku puzzles.
Other colorings
Ramsey theory
An important class of improper coloring problems is studied in Ramsey theory, where the graph's edges are assigned to colors, and there is no restriction on the colors of incident edges. A simple example is the theorem on friends and strangers, which states that in any coloring of the edges of , the complete graph of six vertices, there will be a monochromatic triangle; often illustrated by saying that any group of six people either has three mutual strangers or three mutual acquaintances. Ramsey theory is concerned with generalisations of this idea to seek regularity amid disorder, finding general conditions for the existence of monochromatic subgraphs with given structure.
Other colorings
Adjacent-vertex-distinguishing-total coloring A total coloring with the additional restriction that any two adjacent vertices have different color sets
Acyclic coloring Every 2-chromatic subgraph is acyclic
B-coloring a coloring of the vertices where each color class contains a vertex that has a neighbor in all other color classes.
Circular coloring Motivated by task systems in which production proceeds in a cyclic way
Cocoloring An improper vertex coloring where every color class induces an independent set or a clique
Complete coloring Every pair of colors appears on at least one edge
Defective coloring An improper vertex coloring where every color class induces a bounded degree subgraph.
Distinguishing coloring An improper vertex coloring that destroys all the symmetries of the graph
Equitable coloring The sizes of color classes differ by at most one
Exact coloring Every pair of colors appears on exactly one edge
Fractional coloring Vertices may have multiple colors, and on each edge the sum of the color parts of each vertex is not greater than one
Hamiltonian coloring Uses the length of the longest path between two vertices, also known as the detour distance
Harmonious coloring Every pair of colors appears on at most one edge
Incidence coloring Each adjacent incidence of vertex and edge is colored with distinct colors
Inherited vertex coloring A set of vertex colorings induced by perfect matchings of edge-colored Graphs.
Interval edge coloring A color of edges meeting in a common vertex must be contiguous
List coloring Each vertex chooses from a list of colors
List edge-coloringEach edge chooses from a list of colors
L(h, k)-coloring Difference of colors at adjacent vertices is at least h and difference of colors of vertices at a distance two is at least k. A particular case is L(2,1)-coloring.
Oriented coloring Takes into account orientation of edges of the graph
Path coloring Models a routing problem in graphs
Radio coloring Sum of the distance between the vertices and the difference of their colors is greater than k + 1, where k is a positive integer.
Rank coloring If two vertices have the same color i, then every path between them contain a vertex with color greater than i
Subcoloring An improper vertex coloring where every color class induces a union of cliques
Sum coloring The criterion of minimalization is the sum of colors
Star coloring Every 2-chromatic subgraph is a disjoint collection of stars
Strong coloring Every color appears in every partition of equal size exactly once
Strong edge coloring Edges are colored such that each color class induces a matching (equivalent to coloring the square of the line graph)
T-coloring Absolute value of the difference between two colors of adjacent vertices must not belong to fixed set T
Total coloring Vertices and edges are colored
Centered coloring Every connected induced subgraph has a color that is used exactly once
Triangle-free edge coloring The edges are colored so that each color class forms a triangle-free subgraph
Weak coloring An improper vertex coloring where every non-isolated node has at least one neighbor with a different color
Coloring can also be considered for signed graphs and gain graphs.
| Mathematics | Graph theory | null |
426856 | https://en.wikipedia.org/wiki/Adder%20%28electronics%29 | Adder (electronics) | An adder, or summer, is a digital circuit that performs addition of numbers. In many computers and other kinds of processors, adders are used in the arithmetic logic units (ALUs). They are also used in other parts of the processor, where they are used to calculate addresses, table indices, increment and decrement operators and similar operations.
Although adders can be constructed for many number representations, such as binary-coded decimal or excess-3, the most common adders operate on binary numbers.
In cases where two's complement or ones' complement is being used to represent negative numbers, it is trivial to modify an adder into an adder–subtractor.
Other signed number representations require more logic around the basic adder.
History
George Stibitz invented the 2-bit binary adder (the Model K) in 1937.
Binary adders
Half adder
The half adder adds two single binary digits and . It has two outputs, sum () and carry (). The carry signal represents an overflow into the next digit of a multi-digit addition. The value of the sum is . The simplest half-adder design, pictured on the right, incorporates an XOR gate for and an AND gate for . The Boolean logic for the sum (in this case ) will be whereas for the carry () will be . With the addition of an OR gate to combine their carry outputs, two half adders can be combined to make a full adder. The half adder adds two input bits and generates a carry and sum, which are the two outputs of a half adder. The input variables of a half adder are called the augend and addend bits. The output variables are the sum and carry.
The truth table for the half adder is:
{| class="wikitable" style="text-align:center"
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| A || B || Cout || S
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Various half adder digital logic circuits:
Full adder
A full adder adds binary numbers and accounts for values carried in as well as out. A one-bit full-adder adds three one-bit numbers, often written as , , and ; and are the operands, and is a bit carried in from the previous less-significant stage. The circuit produces a two-bit output. Output carry and sum are typically represented by the signals and , where the sum equals . The full adder is usually a component in a cascade of adders, which add 8, 16, 32, etc. bit binary numbers.
A full adder can be implemented in many different ways such as with a custom transistor-level circuit or composed of other gates. The most common implementation is with:
The above expressions for and can be derived from using a Karnaugh map to simplify the truth table.
In this implementation, the final OR gate before the carry-out output may be replaced by an XOR gate without altering the resulting logic. This is because when A and B are both 1, the term is always 0, and hence can only be 0. Thus, the inputs to the final OR gate can never be both 1's (this is the only combination for which the OR and XOR outputs differ).
Due to the functional completeness property of the NAND and NOR gates, a full adder can also be implemented using nine NAND gates, or nine NOR gates.
Using only two types of gates is convenient if the circuit is being implemented using simple integrated circuit chips which contain only one gate type per chip.
A full adder can also be constructed from two half adders by connecting and to the input of one half adder, then taking its sum-output as one of the inputs to the second half adder and as its other input, and finally the carry outputs from the two half-adders are connected to an OR gate. The sum-output from the second half adder is the final sum output () of the full adder and the output from the OR gate is the final carry output (). The critical path of a full adder runs through both XOR gates and ends at the sum bit . Assumed that an XOR gate takes 1 delays to complete, the delay imposed by the critical path of a full adder is equal to:
The critical path of a carry runs through one XOR gate in adder and through 2 gates (AND and OR) in carry-block and therefore, if AND or OR gates take 1 delay to complete, has a delay of:
The truth table for the full adder is:
{| class="wikitable" style="text-align:center"
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!colspan="3"| Inputs || colspan="2"| Outputs
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| A || B || Cin || Cout || S
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Inverting all inputs of a full adder also inverts all of its outputs, which can be used in the design of fast ripple-carry adders, because there is no need to invert the carry.
Various full adder digital logic circuits:
Adders supporting multiple bits
Ripple-carry adder
It is possible to create a logical circuit using multiple full adders to add N-bit numbers. Each full adder inputs a , which is the of the previous adder. This kind of adder is called a ripple-carry adder (RCA), since each carry bit "ripples" to the next full adder. The first (and only the first) full adder may be replaced by a half adder (under the assumption that ).
The layout of a ripple-carry adder is simple, which allows fast design time; however, the ripple-carry adder is relatively slow, since each full adder must wait for the carry bit to be calculated from the previous full adder. The gate delay can easily be calculated by inspection of the full adder circuit. Each full adder requires three levels of logic. In a 32-bit ripple-carry adder, there are 32 full adders, so the critical path (worst case) delay is 3 (from input to carry in first adder) + 31 × 2 (for carry propagation in latter adders) = 65 gate delays.
The general equation for the worst-case delay for a n-bit carry-ripple adder, accounting for both the sum and carry bits, is:
A design with alternating carry polarities and optimized AND-OR-Invert gates can be about twice as fast.
Carry-lookahead adder (Weinberger and Smith, 1958)
To reduce the computation time, Weinberger and Smith invented a faster way to add two binary numbers by using carry-lookahead adders (CLA). They introduced two signals ( and ) for each bit position, based on whether a carry is propagated through from a less significant bit position (at least one input is a 1), generated in that bit position (both inputs are 1), or killed in that bit position (both inputs are 0). In most cases, is simply the sum output of a half adder and is the carry output of the same adder. After and are generated, the carries for every bit position are created.
Mere derivation of Weinberger-Smith CLA recurrence, are: Brent–Kung adder (BKA), and the Kogge–Stone adder (KSA).
This was shown in Oklobdzija and Zeydel paper in IEEE Journal of Solid-State Circutis.
Some other multi-bit adder architectures break the adder into blocks. It is possible to vary the length of these blocks based on the propagation delay of the circuits to optimize computation time. These block based adders include the carry-skip (or carry-bypass) adder which will determine and values for each block rather than each bit, and the carry-select adder which pre-generates the sum and carry values for either possible carry input (0 or 1) to the block, using multiplexers to select the appropriate result when the carry bit is known.
By combining multiple carry-lookahead adders, even larger adders can be created. This can be used at multiple levels to make even larger adders. For example, the following adder is a 64-bit adder that uses four 16-bit CLAs with two levels of lookahead carry units.
Other adder designs include the carry-select adder, conditional sum adder, carry-skip adder, and carry-complete adder.
Carry-save adders
If an adding circuit is to compute the sum of three or more numbers, it can be advantageous to not propagate the carry result. Instead, three-input adders are used, generating two results: a sum and a carry. The sum and the carry may be fed into two inputs of the subsequent 3-number adder without having to wait for propagation of a carry signal. After all stages of addition, however, a conventional adder (such as the ripple-carry or the lookahead) must be used to combine the final sum and carry results.
3:2 compressors
A full adder can be viewed as a 3:2 lossy compressor: it sums three one-bit inputs and returns the result as a single two-bit number; that is, it maps 8 input values to 4 output values. (the term "compressor" instead of "counter" was introduced in)Thus, for example, a binary input of 101 results in an output of (decimal number 2). The carry-out represents bit one of the result, while the sum represents bit zero. Likewise, a half adder can be used as a 2:2 lossy compressor, compressing four possible inputs into three possible outputs.
Such compressors can be used to speed up the summation of three or more addends. If the number of addends is exactly three, the layout is known as the carry-save adder. If the number of addends is four or more, more than one layer of compressors is necessary, and there are various possible designs for the circuit: the most common are Dadda and Wallace trees. This kind of circuit is most notably used in multiplier circuits, which is why these circuits are also known as Dadda and Wallace multipliers.
Quantum adders
Using only the Toffoli and CNOT quantum logic gates, it is possible to produce quantum full- and half-adders. The same circuits can also be implemented in classical reversible computation, as both CNOT and Toffoli are also classical logic gates.
Since the quantum Fourier transform has a low circuit complexity, it can efficiently be used for adding numbers as well.
Analog adders
Just as in Binary adders, combining two input currents effectively adds those currents together. Within the constraints of the hardware, non-binary signals (i.e. with a base higher than 2) can be added together to calculate a sum. Also known as a "summing amplifier", this technique can be used to reduce the number of transistors in an addition circuit.
| Technology | Digital logic | null |
426889 | https://en.wikipedia.org/wiki/Photometry%20%28optics%29 | Photometry (optics) | Photometry is a branch of optics that deals with measuring light in terms of its perceived brightness to the human eye. It is concerned with quantifying the amount of light that is emitted, transmitted, or received by an object or a system.
In modern photometry, the radiant power at each wavelength is weighted by a luminosity function that models human brightness sensitivity. Typically, this weighting function is the photopic sensitivity function, although the scotopic function or other functions may also be applied in the same way. The weightings are standardized by the CIE and ISO.
Photometry is distinct from radiometry, which is the science of measurement of radiant energy (including light) in terms of absolute power.
Photometry and the eye
The human eye is not equally sensitive to all wavelengths of visible light. Photometry attempts to account for this by weighting the measured power at each wavelength with a factor that represents how sensitive the eye is at that wavelength. The standardized model of the eye's response to light as a function of wavelength is given by the luminosity function. The eye has different responses as a function of wavelength when it is adapted to light conditions (photopic vision) and dark conditions (scotopic vision). Photometry is typically based on the eye's photopic response, and so photometric measurements may not accurately indicate the perceived brightness of sources in dim lighting conditions where colors are not discernible, such as under just moonlight or starlight. Photopic vision is characteristic of the eye's response at luminance levels over three candela per square metre. Scotopic vision occurs below 2 × 10−5 cd/m2. Mesopic vision occurs between these limits and is not well characterised for spectral response.
Photometric quantities
Measurement of the effects of electromagnetic radiation became a field of study as early as the end of the 18th century. Measurement techniques varied depending on the effects under study and gave rise to different nomenclature. The total heating effect of infrared radiation as measured by thermometers led to the development of radiometric units in terms of total energy and power. The use of the human eye as a detector led to photometric units, weighted by the eye's response characteristic. Study of the chemical effects of ultraviolet radiation led to characterization by the total dose or actinometric units expressed in photons per second.
Many different units of measure are used for photometric measurements. The adjective "bright" can refer to a light source which delivers a high luminous flux (measured in lumens), or to a light source which concentrates the luminous flux it has into a very narrow beam (candelas), or to a light source that is seen against a dark background. Because of how light propagates through three-dimensional space — spreading out, becoming concentrated, reflecting off shiny or matte surfaces — and because light consists of many different wavelengths, the number of fundamentally different kinds of light measurements that can be made is large, and so are the numbers of quantities and units that represent them.
For example, offices are typically "brightly" illuminated by an array of many recessed fluorescent lights for a combined high luminous flux. A laser pointer has very low luminous flux (it could not illuminate a room) but is blindingly bright in one direction (high luminous intensity in that direction).
Photometric versus radiometric quantities
There are two parallel systems of quantities known as photometric and radiometric quantities. Every quantity in one system has an analogous quantity in the other system. Some examples of parallel quantities include:
Luminance (photometric) and radiance (radiometric)
Luminous flux (photometric) and radiant flux (radiometric)
Luminous intensity (photometric) and radiant intensity (radiometric)
In photometric quantities every wavelength is weighted according to how sensitive the human eye is to it, while radiometric quantities use unweighted absolute power. For example, the eye responds much more strongly to green light than to red, so a green source will have greater luminous flux than a red source with the same radiant flux would. Radiant energy outside the visible spectrum does not contribute to photometric quantities at all, so for example a 1000 watt space heater may put out a great deal of radiant flux (1000 watts, in fact), but as a light source it puts out very few lumens (because most of the energy is in the infrared, leaving only a dim red glow in the visible).
Watts versus lumens
Watts are units of radiant flux while lumens are units of luminous flux. A comparison of the watt and the lumen illustrates the distinction between radiometric and photometric units.
The watt is a unit of power. We are accustomed to thinking of light bulbs in terms of power in watts. This power is not a measure of the amount of light output, but rather indicates how much energy the bulb will use. Because incandescent bulbs sold for "general service" all have fairly similar characteristics (same spectral power distribution), power consumption provides a rough guide to the light output of incandescent bulbs.
Watts can also be a direct measure of output. In a radiometric sense, an incandescent light bulb is about 80% efficient: 20% of the energy is lost (e.g. by conduction through the lamp base). The remainder is emitted as radiation, mostly in the infrared. Thus, a 60 watt light bulb emits a total radiant flux of about 45 watts. Incandescent bulbs are, in fact, sometimes used as heat sources (as in a chick incubator), but usually they are used for the purpose of providing light. As such, they are very inefficient, because most of the radiant energy they emit is invisible infrared. A compact fluorescent lamp can provide light comparable to a 60 watt incandescent while consuming as little as 15 watts of electricity.
The lumen is the photometric unit of light output. Although most consumers still think of light in terms of power consumed by the bulb, in the U.S. it has been a trade requirement for several decades that light bulb packaging give the output in lumens. The package of a 60 watt incandescent bulb indicates that it provides about 900 lumens, as does the package of the 15 watt compact fluorescent.
The lumen is defined as amount of light given into one steradian by a point source of one candela strength; while the candela, a base SI unit, is defined as the luminous intensity of a source of monochromatic radiation, of frequency 540 terahertz, and a radiant intensity of 1/683 watts per steradian. (540 THz corresponds to about 555 nanometres, the wavelength, in the green, to which the human eye is most sensitive. The number 1/683 was chosen to make the candela about equal to the standard candle, the unit which it superseded).
Combining these definitions, we see that 1/683 watt of 555 nanometre green light provides one lumen.
The relation between watts and lumens is not just a simple scaling factor. We know this already, because the 60 watt incandescent bulb and the 15 watt compact fluorescent can both provide 900 lumens.
The definition tells us that 1 watt of pure green 555 nm light is "worth" 683 lumens. It does not say anything about other wavelengths. Because lumens are photometric units, their relationship to watts depends on the wavelength according to how visible the wavelength is. Infrared and ultraviolet radiation, for example, are invisible and do not count. One watt of infrared radiation (which is where most of the radiation from an incandescent bulb falls) is worth zero lumens. Within the visible spectrum, wavelengths of light are weighted according to a function called the "photopic spectral luminous efficiency." According to this function, 700 nm red light is only about 0.4% as efficient as 555 nm green light. Thus, one watt of 700 nm red light is "worth" only 2.7 lumens.
Because of the summation over the visual portion of the EM spectrum that is part of this weighting, the unit of "lumen" is color-blind: there is no way to tell what color a lumen will appear. This is equivalent to evaluating groceries by number of bags: there is no information about the specific content, just a number that refers to the total weighted quantity.
Photometric measurement techniques
Photometric measurement is based on photodetectors, devices (of several types) that produce an electric signal when exposed to light. Simple applications of this technology include switching luminaires on and off based on ambient light conditions, and light meters, used to measure the total amount of light incident on a point.
More complex forms of photometric measurement are used frequently within the lighting industry. Spherical photometers can be used to measure the directional luminous flux produced by lamps, and consist of a large-diameter globe with a lamp mounted at its center. A photocell rotates about the lamp in three axes, measuring the output of the lamp from all sides.
Lamps and lighting fixtures are tested using goniophotometers and rotating mirror photometers, which keep the photocell stationary at a sufficient distance that the luminaire can be considered a point source. Rotating mirror photometers use a motorized system of mirrors to reflect light emanating from the luminaire in all directions to the distant photocell; goniophotometers use a rotating 2-axis table to change the orientation of the luminaire with respect to the photocell. In either case, luminous intensity is tabulated from this data and used in lighting design.
Non-SI photometry units
Luminance
Footlambert
Millilambert
Stilb
Illuminance
Foot-candle
Phot
| Physical sciences | Optics | Physics |
426982 | https://en.wikipedia.org/wiki/Pedestrian%20zone | Pedestrian zone | Pedestrian zones (also known as auto-free zones and car-free zones, as pedestrian precincts in British English, and as pedestrian malls in the United States and Australia) are areas of a city or town restricted to use by people on foot or human-powered transport such as bicycles, with non-emergency motor traffic not allowed. Converting a street or an area to pedestrian-only use is called pedestrianisation.
Pedestrianisation usually aims to provide better accessibility and mobility for pedestrians, to enhance the amount of shopping and other business activities in the area or to improve the attractiveness of the local environment in terms of aesthetics, air pollution, noise and crashes involving motor vehicle with pedestrians. In some cases, motor traffic in surrounding areas increases, as it is displaced rather than replaced. Nonetheless, pedestrianisation schemes are often associated with significant falls in local air and noise pollution and in accidents, and frequently with increased retail turnover and increased property values locally.
A car-free development generally implies a large-scale pedestrianised area that relies on modes of transport other than the car, while pedestrian zones may vary in size from a single square to entire districts, but with highly variable degrees of dependence on cars for their broader transport links.
Pedestrian zones have a great variety of approaches to human-powered vehicles such as bicycles, inline skates, skateboards and kick scooters. Some have a total ban on anything with wheels, others ban certain categories, others segregate the human-powered wheels from foot traffic, and others still have no rules at all. Many Middle Eastern kasbahs have no motorized traffic, but use donkey- or hand-carts to carry goods.
History
Origins in arcades
The idea of separating pedestrians from wheeled traffic is an old one, dating back at least to the Renaissance. However, the earliest modern implementation of the idea in cities seems to date from about 1800, when the first covered shopping arcade was opened in Paris. Separated shopping arcades were constructed throughout Europe in the 19th century, precursors of modern shopping malls. A number of architects and city planners, including Joseph Paxton, Ebenezer Howard, and Clarence Stein, in the 19th and early 20th centuries proposed plans to separate pedestrians from traffic in various new developments.
1920s–1970s
The first "pedestrianisation" of an existing street seems to have taken place "around 1929" in Essen, Germany. This was in Limbecker Straße, a very narrow shopping street that could not accommodate both vehicular and pedestrian traffic. Two other German cities followed this model in the early 1930s, but the idea was not seen outside Germany. Following the devastation of the Second World War a number of European cities implemented plans to pedestrianise city streets, although usually on a largely ad hoc basis, through the early 1950s, with little landscaping or planning. By 1955 twenty-one German cities had closed at least one street to automobile traffic, although only four were "true" pedestrian streets, designed for the purpose. At this time pedestrianisation was not seen as a traffic restraint policy, but rather as a complement to customers who would arrive by car in a city centre.
Pedestrianisation was also common in the United States during the 1950s and 60s as downtown businesses attempted to compete with new suburban shopping malls. However, most of these initiatives were not successful in the long term, and about 90% have been changed back to motorised areas.
1980s–2010s
In the United States, several pedestrian zones in major tourist areas were successful, such as the renovation of the mall in Santa Monica on Los Angeles' Westside and its relaunch as the Third Street Promenade; the creation of the covered, pedestrian Fremont Street Experience in Downtown Las Vegas; the revival of East 4th Street in Downtown Cleveland; and the new pedestrian zone created in the mid-2010s in New York City including along Broadway (the street) and around Times Square.
COVID-19 pandemic
During the COVID-19 pandemic in 2020, some cities had made the pedestrianization of additional streets to encourage social distancing and in many cases to provide extra rooms for restaurants to serve food on patios extended into the newly available spaces. In the United States, New York City closed up to of streets to cars across the city. In Madrid, Spain, the city pedestrianized of streets and of spaces in total. The COVID-19 pandemic gave also birth to proposals for radical change in the organisation of the city, in particular Barcelona, being the pedestrianisation of the whole city and the proposal of an inversion of the concept of sidewalk two elements of the Manifesto for the Reorganisation of the city, written by architecture and urban theorist Massimo Paolini and signed by 160 academics and 300 architects.
Definitions and types
A pedestrian zone is often limited in scope: for example, a single square or a few streets reserved for pedestrians, within a city where residents still largely get around in cars. A car-free town, city or region may be much larger.
Car free towns, cities and regions
A car-free zone is different from a typical pedestrian zone, in that it implies a development largely predicated on modes of transport other than the car.
Examples
A number of towns and cities in Europe have never allowed motor vehicles. Archetypal examples are:
Venice, which occupies many islands in a lagoon, divided by and accessed from canals. Motor traffic stops at the car park at the head of the viaduct from the mainland, and water transport and walking take over from there. However, motor vehicles are allowed on the nearby Lido.
Zermatt in the Swiss Alps. Most visitors reach Zermatt by a cog railway, and there are pedestrian-only streets, but there are also roads with motor vehicles.
Other examples are:
Cinque Terre in Italy
Ghent in Belgium: the pedestrian zone was extended in 2017 from 35 to more than 50 hectares (123 acres), one of the largest car-free areas in Europe.
Pontevedra in Spain, an international model of pedestrianization, almost 50% of the city is pedestrianised.;
The Old Town of Rhodes, where many, if not most, of the streets are too steep and/or narrow for car traffic.
Mount Athos, an autonomous monastic state under the sovereignty of Greece, does not permit automobiles on its territory. Trucks and work-related vehicles only are in use there.
The medieval city of Mdina in Malta does not allow automobiles past the city walls. It is known as the "Silent City" because of the absence of motor traffic in the city.
Sark, an island in the English Channel, is a car-free zone where only bicycles, carriages and tractors are used as transportation.
Gulangyu, an island off the coast of Xiamen in southeastern China. The only vehicles permitted are small electric buggies and electric government service vehicles.
To assist with transport from the car parks in at the edge of car-free cities, there are often bus stations, bicycle sharing stations, and the like.
Car-free development
The term car-free development implies a physical change: either build-up or changes to an existing built area.
Melia et al. (2010) define car-free developments as "residential or mixed use developments which:
Normally provide a traffic-free immediate environment, and
Offer no parking or limited parking separated from the residence, and:
Are designed to enable residents to live without owning a car."
This definition (which they distinguish from the more common "low car development") is based mainly on experience in North West Europe, where the movement for car-free development began. Within this definition, three types are identified:
Vauban model, based on Vauban, Freiburg: it is not "carfree", but "parking-space-free" () in some streets.
Limited Access model
Pedestrianised centres with residential population
Limited access type
The more common form of carfree development involves some sort of physical barrier, which prevents motor vehicles from penetrating into a car-free interior. Melia et al. describe this as the "limited access" type. In some cases, such as Stellwerk 60 in Cologne, there is a removable barrier, controlled by a residents' organisation. In Amsterdam, Waterwijk is a 6-hectare neighborhood where cars may only access parking areas from the streets that form the edges of the neighborhood; all of the inner areas of the neighborhood are car-free.
Temporary car-free streets
Many cities close certain streets to automobiles, typically on weekends and especially in warm weather, to provide more urban space for recreation, and to increase foot traffic to nearby businesses. Examples include Newbury Street in Boston, and Memorial Drive in Cambridge, Massachusetts (which is along a river). In some cases, popularity has resulted in streets being permanently closed to cars, including JFK Drive in Golden Gate Park, San Francisco; Griffith Drive in Griffith Park, Los Angeles; and Capel Street in Dublin.
Reception
Benefits
Several studies have been carried out on European carfree developments. The most comprehensive was conducted in 2000 by Jan Scheurer. Other more recent studies have been made of specific car-free areas such as Vienna's Floridsdorf car-free development.
Characteristics of car-free developments:
Very low levels of car use, resulting in much less traffic on surrounding roads
High rates of walking and cycling
More independent movement and active play for children
Less land is used for parking and roads, so more available for green or social space
The main benefits found for car-free developments:
Low atmospheric emissions
Low road accident rates
Better built environment conditions
Encouragement of active modes.
The main problems related to parking management. Where parking is not controlled in the surrounding area, this often results in complaints from neighbours about overspill parking.
Problems and criticism
There were calls for traffic to be reinstated in Trafalgar Square, London, after pedestrianization caused noise nuisance for visitors to the National Gallery. The director of the gallery is reported to have blamed pedestrianization for the "trashing of a civic space".
Local shopkeepers may be critical of the effect of pedestrianization on their businesses. Reduced through traffic can lead to fewer customers using local businesses, depending on the environment and the area's dependence on the through traffic.
By region and country
Europe
A large number of European towns and cities have made part of their centres car-free since the early 1960s. These are often accompanied by car parks on the edge of the pedestrianised zone, and, in the larger cases, park and ride schemes.
Armenia
Northern Avenue, located in the Kentron district of central Yerevan, is a large pedestrian avenue. The avenue was inaugurated in 2007 and is mainly home to residential buildings, offices, luxury shops and restaurants.
Belgium
In Belgium, Brussels implemented Europe's largest pedestrian zone (French: Le Piétonnier), in phases starting in 2015 and will cover . The area covers much of the historic center within the Small Ring (the ring road built on the site of the 14th-century walls), including the Grand-Place/Grote Markt, the Place de Brouckère/De Brouckèreplein, the Boulevard Anspach/Anspachlaan, and the Place de la Bourse/Beursplein.
Denmark
Central Copenhagen is one of the oldest and largest: it was converted from car traffic into a pedestrian zone in 1962 as an experiment, and is centered on Strøget, which is not a single street but a series of interconnected avenues which create a very large pedestrian zone, although it is crossed in places by streets with vehicular traffic. Most of these zones allow delivery trucks to service the businesses there during the early morning, and street-cleaning vehicles will usually go through these streets after most shops have closed for the night. It has grown in size from in 1962 to in 1996.
Germany
A number of German islands ban or strictly limit the private use of motor vehicles. Heligoland, Hiddensee, and all but two of the East Frisian islands are car-free; Borkum and Norderney have car-free zones and strictly limit automobile use during the summer season and in certain areas, also forbidding travel at night. Some areas provide exceptions for police and emergency vehicles; Heligoland also bans bicycles.
In the early 1980s, the Alternative Liste für Demokratie und Umweltschutz (which later became part of Alliance 90/The Greens) unsuccessfully campaigned to make West Berlin a car-free zone.
Netherlands
In the Netherlands, the inner city of Arnhem has a pedestrian zone () within the boundaries of the following streets and squares: Nieuwe Plein, Willemsplein, Gele Rijdersplein, Looierstraat, Velperbinnensingel, Koningsplein, St. Catharinaplaats, Beekstraat, Walburgstraat, Turfstraat, Kleine Oord, and Nieuwe Oeverstraat.
Rotterdam's city center was almost completely destroyed by German bombing in May 1940. The city decided to build a central shopping street, for pedestrians only, the Lijnbaan, which became Europe's first purpose-built pedestrian street. The Lijnbaan served as a model for many other such streets in the early post-World War II era, such as Warsaw, Prague, Hamburg, and the UK's first pedestrianised shopping precinct in Stevenage in 1959. Rotterdam has since expanded the pedestrian zone to other streets. As of 2018, Rotterdam featured three different types of pedestrian zones: "pedestrian zones", "pedestrian zones, cycling permitted outside of shopping hours", and "pedestrian zones, cycling permitted 24/7". Three exceptions to motor vehicles could apply to specific sections of these three zones, namely: "logistics allowed within window times (5 to 10:30 a.m)", "logistics allowed 24/7", and "commercial traffic allowed during market days".
United Kingdom
In Britain, shopping streets primarily for pedestrians date back to the thirteenth century. A 1981 study found that many Victorian and later arcades continued to be used. A third of London's 168 precincts at that time had been built before 1939, as were a tenth of the 1,304 precincts in the U.K. as a whole.
Early post-1945 new towns carried on the tradition of providing some traffic-free shopping streets. However, in the conversion of traditional shopping streets to pedestrian precincts, Britain started only in 1967 (versus Germany's first conversion in 1929, or the first in the U.S. in 1959). Since then growth was rapid, such that by 1980 a study found that most British towns and cities had a pedestrian shopping precinct; 1,304 in total.
Turkey
In Istanbul, İstiklal Caddesi is a pedestrian street (except for a historic streetcar that runs along it) and a major tourist draw.
U.S. and Canada
Canada
Some Canadian examples are the Sparks Street Mall area of Ottawa, the Distillery District in Toronto, Scarth Street Mall in Regina, Stephen Avenue Mall in Calgary (with certain areas open to parking for permit holders) and part of Prince Arthur Street and the Gay Village in Montreal. Algonquin and Ward's Islands, parts of the Toronto Islands group, are also car-free zones for all 700 residents. Since summer 2004, Toronto has also been experimenting with "Pedestrian Sundays" in its busy Kensington Market. Granville Mall in Halifax, Nova Scotia was a run-down section of buildings on Granville Street built in the 1840s that was restored in the late 1970s. The area was then closed off to vehicles.
United States
Downtown pedestrian zones
In the United States, these zones are commonly called pedestrian malls or pedestrian streets and today are relatively rare, with a few notable exceptions. In 1959, Kalamazoo was the first American city to implement a "pedestrian mall" in its downtown core. This became a method that some cities applied for their downtowns to compete with the growing suburban shopping malls of the time. In the 1960s and 70s, over 200 towns in the United States adopted this approach.
The Downtown Mall in Charlottesville, VA is one of the longest pedestrian malls in the United States, created in 1976 and spanning nine city blocks.
A number of streets and malls in New York City are now pedestrian-only, including 6½ Avenue, Fulton Street, parts of Broadway, and a block of 25th Street.
A portion of Third Street in Santa Monica in Greater Los Angeles was converted into a pedestrian mall in the 1960s to become what is now the Third Street Promenade, a very popular shopping district located just a few blocks from the beach and Santa Monica Pier.
Lincoln Road in Miami Beach, which had previously been a shopping street with traffic, was converted into a pedestrian only street in 1960. The designer was Morris Lapidus. Lincoln Road Mall is now one of the main attractions in Miami Beach.
The idea of exclusive pedestrian zones lost popularity through the 1980s and into the 1990s and results were generally disappointing, but are enjoying a renaissance with the 1989 renovation and relaunch of the Third Street Promenade in Santa Monica, California, the 1994-5 Fremont Street Experience in Las Vegas and recent pedestrianization of various streets in New York City. These pedestrian zones were more closely tied to the success of retail than in Europe, and by the 1980s, most did not succeed competing with ever more elaborate enclosed malls. Almost all of this generation of pedestrian malls built from 1959 through to the 1970s, have disappeared, or were shrunk down in the 1990s at the request of the retailers. Half of Kalamazoo's pedestrian mall has been converted into a regular street with auto traffic, though with wide sidewalks.
Outside large cities
Mackinac Island, between the upper and lower peninsulas of Michigan, banned horseless carriages in 1896, making it auto-free. The original ban still stands, except for emergency vehicles. Travel on the island is largely by foot, bicycle, or horse-drawn carriage. An road, M-185 rings the island, and numerous roads cover the interior. M-185 is the only highway in the United States without motorized vehicles.
Fire Island in Suffolk County, New York is pedestrianised east of the Fire Island Lighthouse and west of Smith Point County Park (with the exception of emergency vehicles).
Supai, Arizona, located within the Havasupai Indian Reservation is entirely car-free, the only community in the United States where mail is still carried out by mule. Supai is located eight miles from the nearest road, and is accessible only by foot, horse/mule, or helicopter.
Culdesac Tempe, a 17-acre (0.069 square kilometers) car-free district in Tempe, Arizona, is intended to be the nation's first market-rate rental apartment district to ban its tenants from owning cars. Bikes and emergency vehicles are allowed. It has received significant investments from executives at Lyft and Opendoor.
Latin America
Argentina
Argentina's big cities, Córdoba, Mendoza and Rosario, have lively pedestrianised street centers () combined with town squares and parks which are crowded with people walking at every hour of the day and night.
In Buenos Aires, some stretches of Calle Florida have been pedestrianised since 1913, which makes it one of the oldest car-free thoroughfares in the world today. Pedestrianised Florida, Lavalle and other streets contribute to a vibrant shopping and restaurant scene where street performers and tango dancers abound, streets are crossed with vehicular traffic at chamfered corners.
Brazil
Paquetá Island in Rio de Janeiro is auto-free. The only cars allowed on the island are police and ambulance vehicles. In Rio de Janeiro, the roads beside the beaches are auto-free on Sundays and holidays.
Downtown Rio de Janeiro, Ouvidor Street, over almost its entire length, has been continually a pedestrian space since the mid-nineteenth century when not even carts or carriages were allowed. And the Saara District, also downtown, consists of some dozen or more blocks of colonial streets, off-limits to cars, and crowded with daytime shoppers. Likewise, many of the city's hillside favelas are effectively pedestrian zones as the streets are too narrow and/or steep for automobiles.
Eixo Rodoviário, in Brasília, which is 13 kilometers long and 30 meters wide and is an arterial road connecting the center of that city from both southward and northward wings of Brasília, perpendicular to the well known Eixo Monumental (Monumental Axis in English), is auto-free on Sundays and holidays.
Rua XV de Novembro (15 November Street) in Curitiba is one of the first major pedestrian streets in Brazil.
Chile
Chile has many large pedestrian streets. An example is Paseo Ahumada and Paseo Estado in Santiago, Paseo Barros Arana in Concepción and Calle Valparaíso in Viña del Mar.
Colombia
During his 1998–2001 term, the former Bogotá mayor, U.S.-born Enrique Peñalosa, created several pedestrian streets, plazas and bike paths integrated with a new bus rapid transit system.
The historic center of Cartagena closes some streets to cars during certain hours.
In downtown Armenia, Colombia there is a large pedestrian street where several boutiques are located.
Santa Marta also has permanent pedestrian zones in the historic center around the Cathedral Basílica of Santa Marta.
Mexico
The Historic center of Mexico City has 12 pedestrian streets including Madero Street, and as of 30 June 2020, is expanding the number to 42 pedestrian streets.
Génova is a busy pedestrian street in the Zona Rosa as is Plaza Garibaldi downtown, where mariachis play.
The old city of Guanajuato is largely pedestrian. The steep and/or narrow side streets were never accessible by cars and most other streets were pedestrianized in the 1960s after through traffic was moved to a system of former flood control tunnels that was no longer necessary due to a new dam.
Playa del Carmen has a pedestrian mall, Quinta Avenida, ("Fifth Avenue") that stretches and receives 4 million visitors annually with hundreds of shops and restaurants.
Peru
Jirón de La Unión in Lima is a traditional pedestrian street located in the Historic Centre of Lima, part of the capital of Peru.
In the city of Arequipa, Mercaderes is also a considerably large pedestrian street.
Also, recently three of the four streets surrounding the city's main square or "Plaza de Armas" were also made pedestrian.
South and East Asia
Mainland China
Nanjing Road in Shanghai is perhaps the most well-known pedestrian zone in mainland China. Wangfujing is a famous tourist and retail oriented pedestrian zone in Beijing. Chunxilu in Chengdu is the most well known in western China. Dongmen is the busiest business zone in Shenzhen. Zhongyang Street is a historical large pedestrian street in Harbin.
Hong Kong
In Hong Kong, since 2000, the government has been implementing full-time or part-time pedestrian streets in a number of areas, including Causeway Bay, Central, Wan Chai, Mong Kok, and Tsim Sha Tsui. The most popular pedestrian street is Sai Yeung Choi Street. It was converted into a pedestrian street in 2003. From December 2008 to May 2009, there were three acid attacks during which corrosive liquids were placed in plastic bottles and thrown from the roof of apartments down onto the street.
India
Vehicles have been banned in the town of Matheran, in Maharashtra, India since the time it was discovered in 1854.
In India, a citizens' initiative in Goa state, has made 18 June Road, Panjim's main shopping boulevard a Non-Motorised Zone(NoMoZo). The road is converted into a NoMoZo for half a day on one Sunday every month.
In Pune, Maharashtra, similar efforts have been made to convert M.G. Road (a.k.a. Main Street) into an open-air mall. The project in question aimed to create a so-called "Walking Plaza".
In May 2019, the North Delhi Municipal Corporation (NDMC) made the busy Ajmal Khan Road in Karol Bagh pedestrian-only.
Church Street in Bangalore went through a pedestrianization process
Japan
Pedestrian zones in Japan are called hokōsha tengoku (歩行者天国, literally "pedestrian heaven"). Clis Road, in Sendai, Japan, is a covered pedestrian mall, as is Hondōri in Hiroshima. Several major streets in Tokyo are closed to vehicles during weekends. One particular temporary hokōsha tengoku in Akihabara was cancelled after the Akihabara massacre in which a man rammed a truck into the pedestrian traffic and subsequently stabbed more than 12 people.
South Korea
Insadong in Seoul, South Korea has a large pedestrian zone (Insadong-gil) during certain hours.
Also in South Korea, in 2013, in the Haenggun-dong neighbourhood of Suwon, streets were closed to cars as a month-long car-free experiment while the city hosted the EcoMobility World Festival. Instead of cars, residents used non-motorized vehicles provided by the festival organizers. The experiment was not unopposed; however, on balance it was considered a success. Following the festival, the city embarked on discussions about adopting the practice on a permanent basis.
Taiwan
Ximending in Taipei, Taiwan is a neighbourhood and shopping district in the Wanhua District of Taipei, Taiwan. It was the first pedestrian zone in Taiwan. The district is very popular in Taiwan. In central Taiwan, Yizhong Street is one of the most popular pedestrian shopping area in Taichung. In Southern Taiwan, the most famous pedestrian shopping area is Shinkuchan in Kaohsiung.
Thailand
In Thailand, some small streets (soi) in Bangkok are designed to be all-time closed to automobile traffic, the city's famous shopping streets of Sampheng Lane in Chinatown and Wang Lang Market nearby to Siriraj Hospital, are the most popular for both local and tourists shopping streets. Additionally the city has built long skywalk systems. Walking Street, Pattaya is also closed to auto traffic. Night markets are routinely closed to auto traffic.
Vietnam
Huế in Vietnam has made 3 roads into pedestrians-only on weekend nights. Also, Hanoi has opened an Old Quarter Walking Street on weekend nights.
Ho Chi Minh City also changed Nguyễn Huệ street into pedestrian zone.
Middle East and North Africa
North Africa contains some of the largest auto-free areas in the world. Fes-al-Bali, a medina of Fes, Morocco, with its population of 156,000, may be the world's largest contiguous completely carfree area, and the medinas of Cairo, Tunis, Casablanca, Meknes, Essaouira, and Tangier are quite extensive.
In Israel, Tel Aviv has a pedestrian mall, near Nahalat Binyamin Street. Ben Yehuda Street in Jerusalem is a pedestrian mall.
Oceania
Australia
In Australia, as in the US, these zones are commonly called pedestrian malls and in most cases comprise only one street. Most pedestrian streets were created in the late 1970s and 1980s, the first being City Walk, Garema Place in Canberra in 1971. Of 58 pedestrian streets created in Australia in the last quarter of the 20th century, 48 remain today, ten having re-introduced car access between 1990 and 2004. All capital cities in Australia have at least one pedestrian street of which most central are: George Street, Pitt Street Mall and Martin Place in Sydney, Bourke Street Mall in Melbourne, Queen Street Mall and Brunswick Street Mall in Brisbane, Rundle Mall in Adelaide, Hay Street and Murray Street Malls in Perth, Elizabeth Street Mall in Hobart, City Walk in Canberra, and Smith Street in Darwin. Many other mid-sized and regional Australian cities also feature pedestrian malls, examples include Rooke Street Devonport Langtree Avenue Mildura, Cavill Avenue Surfers Paradise, Bridge Street Ballarat, Nicholas Street Ipswich, Hargreaves Street Bendigo, Maude Street Shepparton and Little Mallop Street Geelong.
Empirical studies by Jan Gehl indicate an increase of pedestrian traffic as result of public domain improvements in the centres of Melbourne with 39% increase between 1994 and 2004 and Perth with 13% increase between 1993 and 2009.
Most intensive pedestrian traffic flows on a summer weekday have been recorded in Bourke Street Mall Melbourne with 81,000 pedestrians (2004), Rundle Mall Adelaide with 61,360 pedestrians (2002), Pitt Street Mall Sydney with 58,140 (2007) and Murray Street Mall Perth with 48,350 pedestrians (2009).
Rottnest Island off Perth is car-free, only allowing vehicles for essential services. Bicycles are the main form of transport on the island; they can be hired or brought over on the ferry.
In Melbourne's north-eastern suburbs, there have been many proposals to make the Doncaster Hill development area a pedestrian zone. If the proposals are passed, the zone could be one of the largest in the world, by area.
New Zealand
Wellington's Cuba Street became the first pedestrian-only street in New Zealand when in 1965 the Wellington tramway lines were removed and the street was closed off to auto traffic, and after public pressure to keep it closed to automobiles, part of the street was pedestrianised in 1969 and reopened as Cuba Mall.
New Zealand's second-largest city, Christchurch, made its main shopping streets (Cashel & High Street) pedestrian-only in 1982 and created City Mall, also commonly known as Cashel Mall. The concept was first proposed in 1965, around the same time Wellington proposed Cuba Street's pedestrianisation. After the success of Cuba Mall in Wellington, Christchurch decided to continue with the plans. In 1976 the Bridge of Remembrance was pedestrianised, and eventually in August 1982 the entire City Mall was pedestrianised and fully opened to the public. The area was repaved in the late 2000s and again after the Christchurch Earthquakes in 2010 & 2011.
Queenstown has made most of its town centre a pedestrian zone with the lower part of Ballarat Street converted in the 1970s and turned into Queenstown Mall. Most recently, Lower Beach Street has been partially pedestrianised with now only one-way traffic for cars.
Auckland's Lower Queen Street was pedestrianised in 2020.
Town Centre–style pedestrian malls rose in popularity in the 1970 & 1980s, springing up around New Zealand after the success of Cuba Mall. Many, however, have since fallen into disrepair and abandonment and are now classified as Dead malls, including Bishopdale Village Mall, Otara Town Centre, and New Brighton Mall. Pedestrian malls are still being built, however much more scarcely and now are usually called Town Centres and have parking on the outskirts, including Rolleston Fields, The Sands Town Centre, and The Landing Wigram.
A proposal has been made for a pedestrian priority community near Papakura in Auckland. The community would be called Sunfield and cost $4 Billion NZD to build. It is designed to have 4,400 homes and is projected to decrease normal car usage by 90% compared to typical suburbs. It has run into challenges after the project being rejected by Kāinga Ora for fast-tracking following Covid-19; construction authorities took Kāinga Ora to court over the matter.
| Technology | Road infrastructure | null |
427118 | https://en.wikipedia.org/wiki/Principle%20of%20locality | Principle of locality | In physics, the principle of locality states that an object is influenced directly only by its immediate surroundings. A theory that includes the principle of locality is said to be a "local theory". This is an alternative to the concept of instantaneous, or "non-local" action at a distance. Locality evolved out of the field theories of classical physics. The idea is that for a cause at one point to have an effect at another point, something in the space between those points must mediate the action. To exert an influence, something, such as a wave or particle, must travel through the space between the two points, carrying the influence.
The special theory of relativity limits the maximum speed at which causal influence can travel to the speed of light, . Therefore, the principle of locality implies that an event at one point cannot cause a truly simultaneous result at another point. An event at point cannot cause a result at point in a time less than , where is the distance between the points and is the speed of light in vacuum.
The principle of locality plays a critical role in one of the central results of quantum mechanics. In 1935, Albert Einstein, Boris Podolsky, and Nathan Rosen, with their EPR paradox thought experiment, raised the possibility that quantum mechanics might not be a complete theory. They described two systems physically separated after interacting; this pair would be called entangled in modern terminology. They reasoned that without additions, now called hidden variables, quantum mechanics would predict illogical relationships between the physically separated measurements.
In 1964, John Stewart Bell formulated Bell's theorem, an inequality which, if violated in actual experiments, implies that quantum mechanics violates local causality (referred to as local realism in later work), a result now considered equivalent to precluding local hidden variables. Progressive variations on those Bell test experiments have since shown that quantum mechanics broadly violates Bell's inequalities. According to some interpretations of quantum mechanics, this result implies that some quantum effects violate the principle of locality.
Pre-quantum mechanics
During the 17th century, Newton's principle of universal gravitation was formulated in terms of "action at a distance", thereby violating the principle of locality. Newton himself considered this violation to be absurd:
Coulomb's law of electric forces was initially also formulated as instantaneous action at a distance, but in 1880, James Clerk Maxwell showed that field equations – which obey locality – predict all of the phenomena of electromagnetism. These equations show that electromagnetic forces propagate at the speed of light.
In 1905, Albert Einstein's special theory of relativity postulated that no matter or energy can travel faster than the speed of light, and Einstein thereby sought to reformulate physics in a way that obeyed the principle of locality. He later succeeded in producing an alternative theory of gravitation, general relativity, which obeys the principle of locality.
However, a different challenge to the principle of locality developed subsequently from the theory of quantum mechanics, which Einstein himself had helped to create.
Models for locality
Simple spacetime diagrams can help clarify the issues related to locality. A way to describe the issues of locality suitable for discussion of quantum mechanics is illustrated in the diagram. A particle is created in one location, then split and measured in two other, spatially separated, locations. The two measurements are named for Alice and Bob. Alice performs measurements (A) and gets a result ); Bob performs () and gets result . The experiment is repeated many times and the results are compared.
Alice and Bob in spacetime
A spacetime diagram has a time coordinate going vertical and a space coordinate going horizontal. Alice, in a local region on the left, can affect events only in a cone extending in the future as shown; the finite speed of light prevent her from affecting other areas including Bob's location in this case. Similarly we can use the diagram to reason that Bob's local circumstances cannot be altered by Alice at the same time: all events that cause an effect on Bob are in the cone below his location on the diagram. Dashed lines around Alice show her valid future locations; dashed lines around Bob show events that could have caused his present circumstance. When Alice measures quantum states in her location she gets the results labeled ; similarly Bob gets . Models of locality attempt to explain the statistical relationship between these measured values.
Action at a distance
The simplest locality model is no locality: instantaneous action at a distance with no limits for relativity. The locality model for action at a distance is called continuous action. The gray area (a circle here) is a mathematical concept called a "screen". Any path from a location through the screen becomes part of the physical model at that location. The gray ring indicates events from all parts of space and time can affect the probability measured by Alice or Bob. So in the case of continuous action, events at all times and places affect Alice's and Bob's model. This simple model is highly successful for solar planetary dynamics with Newtonian gravity and in electrostatics, cases where relativistic effects are insignificant.
No future-input dependence
Many locality models explicitly or implicitly ignore the possible effect of future events. The spacetime diagram at the right shows the effect of such a restriction when combined with continuous action. Inputs from the future (above the dashed line) are no longer considered part of Alice's or Bob's model. Comparing this diagram with the one for continuous action makes it clear that these are not the same locality model. Common sense arguments about the future not affecting the present are reasonable criteria but such assumptions alter the mathematical character of the models.
Bell's local causality
John Stewart Bell when discussing his Bell's theorem uses the screening model shown at the right. Events in the common past of Alice and Bob are part of the model used in calculating probabilities for Alice and for Bob as indicated the way the screen absorbs those events. However events at Bob's location during Alice measurement and events in the future are excluded. Bell called this assumption local causality, but with the diagram we can reason about the meaning of the assumption without getting tripped up by other meanings of local combined with other meanings of causal. Dash lines show relativistically valid regions in the past of Alice or Bob. The gray arc is the assumed Bell "screen".
Quantum mechanics
The relative positions of our few, easily distinguishable planets (for example) can be seen directly: understanding and measuring their relative location poses only technical issues. The submicroscopic world on the other hand is known only by measurements that average over many seemingly random ("statistical" or "probabilistic") events and measurements can show either particle-like or wave-like results depending on their design. This world is governed by quantum mechanics. The concepts of locality are more complex and they are described in the language of probability and correlation.
In the 1935 Einstein–Podolsky–Rosen paradox paper (EPR paper), Albert Einstein, Boris Podolsky and Nathan Rosen imagined such an experiment. They observed that quantum mechanics predicts what is now known as quantum entanglement and examined its consequences. In their view, the classical principle of locality implied that "no real change can take place" at Bob's site as a result of whatever measurements Alice was doing. Since quantum mechanics does predict a wavefunction collapse that depends on Alice's choice of measurement, they concluded that this was a form of action-at-distance and that the wavefunction could not be a complete description of reality. Other physicists did not agree: they accepted the quantum wavefunction as complete and questioned the nature of locality and reality assumed in the EPR paper.
In 1964 John Stewart Bell investigated whether it might be possible to fulfill Einstein's goal—to "complete" quantum theory—with local hidden variables to explain the correlations between spatially separated particles as predicted by quantum theory. Bell established a criterion to distinguish between local hidden-variables theory and quantum theory by measuring specific values of correlations between entangled particles. Subsequent experimental tests have shown that some quantum effects do violate Bell's inequalities and cannot be reproduced by a local hidden-variables theory. Bell's theorem depends on careful defined models of locality.
Locality and hidden variables
Bell described local causality in terms of probability needed for analysis of quantum mechanics. Using the notation that for the probability of a result with given state , Bell investigated the probability distribution
where represents hidden state variables set (locally) when the two particles are initially co-located. If local causality holds, then the probabilities observed by Alice and by Bob should be only coupled by the hidden variables, and we can show that
Bell proved that a consequence of this factorization are limits on the correlations observed by Alice and Bob known as Bell inequalities. Since quantum mechanics predicts correlations stronger than this limit, locally set hidden variables cannot be added to "complete" quantum theory as desired by the EPR paper.
Numerous experiments specifically designed to probe the issues of locality confirm the predictions of quantum mechanics; these include experiments where the two measurement locations are more than a kilometer apart.
The 2022 Nobel Prize in Physics was awarded to Alain Aspect, John Clauser and Anton Zeilinger, in part "for experiments with entangled photons, establishing the violation of Bell inequalities". The specific aspect of quantum theory that leads to these correlations is termed quantum entanglement, and versions of Bell's scenario are now used to verify entanglement experimentally.
Terminology
Bell's mathematical results, when compared to experimental data, eliminate local hidden-variable mathematical quantum theories. But the interpretation of the math with respect to the physical world remains under debate. Bell described the assumptions behind his work as "local causality", shortened to "locality"; later authors referred to the assumptions as local realism. These different names do not alter the mathematical assumptions.
A review of papers using this phrase suggests that a common (classical) physics definition of realism is
This definition includes classical concepts like "well-defined", which conflicts with quantum superposition, and "prior to ... measurements", which implies (metaphysical) preexistence of properties. Specifically, the term local realism in the context of Bell's theorem cannot be viewed as a kind of "realism" involving locality other than the kind implied by the Bell screening assumption. This conflict between common ideas of realism and quantum mechanics requires careful analysis whenever local realism is discussed.
Adding a "locality" modifier, that the results of two spatially well-separated measurements cannot causally affect each other, does not make the combination relate to Bell's proof; the only interpretation that Bell assumed was the one he called local causality. Consequently, Bell's theorem does not restrict the possibility of nonlocal variables as well as theories based on retrocausality or superdeterminism.
Because of the probabilistic nature of wave function collapse, this apparent violation of locality in quantum mechanics cannot be used to transmit information faster than light, in accordance to the no communication theorem. Asher Peres distinguishes between weak and strong nonlocality, the latter referring to the theories that allow faster-than-light communication. Under these terms, quantum mechanics would allow weakly nonlocal correlations but not strong nonlocality.
Relativistic quantum mechanics
One of the main principles of quantum field theory is the principle of locality. The field operators and the Lagrangian density describing the dynamics of the fields are local, in the sense that interactions are not described by action-at-a-distance. This condition can be achieved by avoiding terms in the Lagrangian that are products of two fields that depend on distant coordinates. Specifically, in relativistic quantum field theory, to enforce the principles of locality and causality the following condition is required: if there are two observables, each localized within two distinct spacetime regions which happen to be at a spacelike separation from each other, the observables must commute. This condition is sometimes imposed as one of the axioms of relativistic quantum field theory.
| Physical sciences | Physics basics: General | Physics |
427445 | https://en.wikipedia.org/wiki/Hoof | Hoof | The hoof (: hooves) is the tip of a toe of an ungulate mammal, which is covered and strengthened with a thick and horny keratin covering. Artiodactyls are even-toed ungulates, species whose feet have an even number of digits; the ruminants with two digits are the most numerous, e.g. giraffe, deer, bison, cattle, goat, pigs, and sheep. The feet of perissodactyl mammals have an odd number of toes, e.g. the horse, the rhinoceros, and the tapir. Although hooves are limb structures primarily found in placental mammals, hadrosaurs such as Edmontosaurus possessed hoofed forelimbs. The marsupial Chaeropus also had hooves.
Description
The hoof surrounds the distal end of the second phalanx, the distal phalanx, and the navicular bone. The hoof consists of the hoof wall, the bars of the hoof, the sole and frog and soft tissue shock absorption structures. The weight of the animal is normally borne by both the sole and the edge of the hoof wall. Hooves perform many functions, including supporting the weight of the animal, dissipating the energy impact as the hooves strike the ground or surface, protecting the tissues and bone within the hoof capsule, and providing traction for the animal. Numerous factors can affect hoof structure and health, including genetics, hoof conformation, environmental influences, and athletic performance of the animal. The ideal hoof has a parallel hoof-pastern axis, a thick hoof wall, adequate sole depth, a solid heel base and growth rings of equal size under the coronary band.
There are four layers within the exterior wall of the hoof. From the outside, a hoof is made up of the stratum externum, the stratum medium, the stratum internum and the dermis parietis. The stratum externum and the stratum medium are difficult to distinguish, the stratum externum is thin and the stratum medium is what makes up the bulk of the hoof wall. Inside the hoof wall is a laminar junction, a soft tissue structure that allows the hoof to withstand the demands of force transmission it undergoes. This tissue structure binds the inner surface of the hoof wall, the dermis parietis and the outer surface of the third phalanx.
Most even-toed ungulates (such as sheep, goats, deer, cattle, bison and pigs) have two main hooves on each foot, together called a cloven hoof. Most of these cloven-hooved animals also have two smaller hooves called dewclaws a little further up the leg – these are not normally used for walking, but in some species with larger dewclaws (such as deer and pigs) they may touch the ground when running or jumping, or if the ground is soft. In the mountain goat, the dewclaw serves to provide extra traction when descending rocky slopes as well as additional drag on loose or slippery surfaces made of ice, dirt, or snow. Other cloven-hooved animals (such as giraffes and pronghorns) have no dewclaws.
In some so-called "cloven-hooved" animals, such as camels, the "hoof" is not properly a hoof – it is not a hard or rubbery sole with a hard wall formed by a thick nail – instead it is a soft toe with little more than a nail merely having an appearance of a hoof.
Some odd-toed ungulates (equids) have one hoof on each foot; others have (or had) three distinct hooved or heavily nailed toes, or one hoof and two dewclaws. The tapir is a special case, having three toes on each hind foot and four toes on each front foot.
Management
Hooves grow continuously. In nature, wild animals are capable of wearing down the hoof as it continuously grows, but captive domesticated species often must undergo specific hoof care for a healthy, functional hoof. Proper care improves biomechanical efficiency and prevents lameness. If not worn down enough by use, such as in the dairy industry, hooves may need to be trimmed. However, too much wear can result in damage of the hooves, and for this reason, horseshoes and oxshoes are used by animals that routinely walk on hard surfaces and carry heavy weight.
Horses
Within the equine world, the expression, "no foot, no horse" emphasizes the importance of hoof health. Hoof care is important in the equine industry. Problems that can arise with poor horse hoof care include hoof cracks, thrush, abscesses and laminitis.
Cattle
A cow hoof is cloven, or divided, into two approximately equal parts, usually called claws. Approximately 95% of lameness in dairy cattle occurs in the feet. Lameness in dairy cows can reduce milk production and fertility, and cause reproductive problems and suffering. For dairy farm profitability, lameness, behind only infertility and mastitis, is the third most important cow health issue.
Hoof trimmers trim and care for bovine hooves, usually dairy cows. Hooves can be trimmed with a sharp knife while the cow is restrained and positioned with ropes. Professional hoof-trimming tend to use angle grinders and some type of hoof trimming crush to make the process quicker and less physically demanding on the hoof trimmer. A hoof trimmer using modern machinery may trim the hooves of more than 10,000 cows per year. The trimmer shapes the hooves to provide the optimal weight-bearing surface. A freshly trimmed hoof may be treated with copper sulfate pentahydrate to prevent foot rot.
Gallery
In culture
Hooves have historical significance in ceremonies and games. They have been used in burial ceremonies.
| Biology and health sciences | External anatomy and regions of the body | Biology |
427499 | https://en.wikipedia.org/wiki/Proteinogenic%20amino%20acid | Proteinogenic amino acid | Proteinogenic amino acids are amino acids that are incorporated biosynthetically into proteins during translation from RNA. The word "proteinogenic" means "protein creating". Throughout known life, there are 22 genetically encoded (proteinogenic) amino acids, 20 in the standard genetic code and an additional 2 (selenocysteine and pyrrolysine) that can be incorporated by special translation mechanisms.
In contrast, non-proteinogenic amino acids are amino acids that are either not incorporated into proteins (like GABA, L-DOPA, or triiodothyronine), misincorporated in place of a genetically encoded amino acid, or not produced directly and in isolation by standard cellular machinery (like hydroxyproline). The latter often results from post-translational modification of proteins. Some non-proteinogenic amino acids are incorporated into nonribosomal peptides which are synthesized by non-ribosomal peptide synthetases.
Both eukaryotes and prokaryotes can incorporate selenocysteine into their proteins via a nucleotide sequence known as a SECIS element, which directs the cell to translate a nearby UGA codon as selenocysteine (UGA is normally a stop codon). In some methanogenic prokaryotes, the UAG codon (normally a stop codon) can also be translated to pyrrolysine.
In eukaryotes, there are only 21 proteinogenic amino acids, the 20 of the standard genetic code, plus selenocysteine. Humans can synthesize 12 of these from each other or from other molecules of intermediary metabolism. The other nine must be consumed (usually as their protein derivatives), and so they are called essential amino acids. The essential amino acids are histidine, isoleucine, leucine, lysine, methionine, phenylalanine, threonine, tryptophan, and valine (i.e. H, I, L, K, M, F, T, W, V).
The proteinogenic amino acids have been found to be related to the set of amino acids that can be recognized by ribozyme autoaminoacylation systems. Thus, non-proteinogenic amino acids would have been excluded by the contingent evolutionary success of nucleotide-based life forms. Other reasons have been offered to explain why certain specific non-proteinogenic amino acids are not generally incorporated into proteins; for example, ornithine and homoserine cyclize against the peptide backbone and fragment the protein with relatively short half-lives, while others are toxic because they can be mistakenly incorporated into proteins, such as the arginine analog canavanine.
The evolutionary selection of certain proteinogenic amino acids from the primordial soup has been suggested to be because of their better incorporation into a polypeptide chain as opposed to non-proteinogenic amino acids.
Structures
The following illustrates the structures and abbreviations of the 21 amino acids that are directly encoded for protein synthesis by the genetic code of eukaryotes. The structures given below are standard chemical structures, not the typical zwitterion forms that exist in aqueous solutions.
IUPAC/IUBMB now also recommends standard abbreviations for the following two amino acids:
Chemical properties
Following is a table listing the one-letter symbols, the three-letter symbols, and the chemical properties of the side chains of the standard amino acids. The masses listed are based on weighted averages of the elemental isotopes at their natural abundances. Forming a peptide bond results in elimination of a molecule of water. Therefore, the protein's mass is equal to the mass of amino acids the protein is composed of minus 18.01524 Da per peptide bond.
General chemical properties
Side-chain properties
§: Values for Asp, Cys, Glu, His, Lys & Tyr were determined using the amino acid residue placed centrally in an alanine pentapeptide. The value for Arg is from Pace et al. (2009). The value for Sec is from Byun & Kang (2011).
N.D.: The pKa value of Pyrrolysine has not been reported.
Note: The pKa value of an amino-acid residue in a small peptide is typically slightly different when it is inside a protein. Protein pKa calculations are sometimes used to calculate the change in the pKa value of an amino-acid residue in this situation.
Gene expression and biochemistry
* UAG is normally the amber stop codon, but in organisms containing the biological machinery encoded by the pylTSBCD cluster of genes the amino acid pyrrolysine will be incorporated.
** UGA is normally the opal (or umber) stop codon, but encodes selenocysteine if a SECIS element is present.
† The stop codon is not an amino acid, but is included for completeness.
†† UAG and UGA do not always act as stop codons (see above).
‡ An essential amino acid cannot be synthesized in humans and must, therefore, be supplied in the diet. Conditionally essential amino acids are not normally required in the diet, but must be supplied exogenously to specific populations that do not synthesize it in adequate amounts.
& Occurrence of amino acids is based on 135 Archaea, 3775 Bacteria, 614 Eukaryota proteomes and human proteome (21 006 proteins) respectively.
Mass spectrometry
In mass spectrometry of peptides and proteins, knowledge of the masses of the residues is useful. The mass of the peptide or protein is the sum of the residue masses plus the mass of water (Monoisotopic mass = 18.01056 Da; average mass = 18.0153 Da). The residue masses are calculated from the tabulated chemical formulas and atomic weights. In mass spectrometry, ions may also include one or more protons (Monoisotopic mass = 1.00728 Da; average mass* = 1.0074 Da). *Protons cannot have an average mass, this confusingly infers to Deuterons as a valid isotope, but they should be a different species (see Hydron (chemistry))
§ Monoisotopic mass
Stoichiometry and metabolic cost in cell
The table below lists the abundance of amino acids in E.coli cells and the metabolic cost (ATP) for synthesis of the amino acids. Negative numbers indicate the metabolic processes are energy favorable and do not cost net ATP of the cell. The abundance of amino acids includes amino acids in free form and in polymerization form (proteins).
Remarks
Catabolism
Amino acids can be classified according to the properties of their main products:
Glucogenic, with the products having the ability to form glucose by gluconeogenesis
Ketogenic, with the products not having the ability to form glucose: These products may still be used for ketogenesis or lipid synthesis.
Amino acids catabolized into both glucogenic and ketogenic products
| Biology and health sciences | Amino acids | Biology |
427655 | https://en.wikipedia.org/wiki/Boletus%20edulis | Boletus edulis | Boletus edulis (English: cep, penny bun, porcino or porcini) is a basidiomycete fungus, and the type species of the genus Boletus.
Prized as an ingredient in various culinary dishes, B. edulis is an edible mushroom held in high regard in many cuisines, and is commonly prepared and eaten in soups, pasta, or risotto. The mushroom is low in fat and digestible carbohydrates, and high in protein, vitamins, minerals and dietary fibre. Although it is sold commercially, it is very difficult to cultivate. Available fresh in autumn throughout Europe and Russia, it is most often dried, packaged, and distributed worldwide. It keeps its flavour after drying, and it is then reconstituted and used in cooking. B. edulis is also one of the few fungi sold pickled.
The fungus grows in deciduous and coniferous forests and tree plantations, forming symbiotic ectomycorrhizal associations with living trees by enveloping the tree's underground roots with sheaths of fungal tissue. The fungus produces spore-bearing fruit bodies above ground in summer and autumn. The fruit body has a large brown cap which on occasion can reach , rarely in diameter and in weight. Like other boletes, it has tubes extending downward from the underside of the cap, rather than gills; spores escape at maturity through the tube openings, or pores. The pore surface of the B. edulis fruit body is whitish when young, but ages to a greenish-yellow. The stout stipe, or stem, is white or yellowish in colour, up to , rarely tall and thick, and partially covered with a raised network pattern, or reticulations.
Widely distributed in the Northern Hemisphere across Europe, Asia, and North America, it does not occur naturally in the Southern Hemisphere, although it has been introduced to southern Africa, Australia, New Zealand, and Brazil. Several closely related European mushrooms formerly thought to be varieties or forms of B. edulis have been shown using molecular phylogenetic analysis to be distinct species, and others previously classed as separate species are conspecific with this species. The western North American species commonly known as the California king bolete (Boletus edulis var. grandedulis) is a large, darker-coloured variant first formally identified in 2007.
Taxonomy
Boletus edulis was first described in 1782 by the French botanist Pierre Bulliard and still bears its original name. The starting date of fungal taxonomy had been set as January 1, 1821, to coincide with the date of the works of the 'father of mycology', Swedish naturalist Elias Magnus Fries, which meant the name required sanction by Fries (indicated in the name by a colon) to be considered valid, as Bulliard's work preceded this date. It was thus written Boletus edulis Bull.:Fr. A 1987 revision of the International Code of Botanical Nomenclature set the starting date at May 1, 1753, the date of publication of Linnaeus' work, the Species Plantarum. Hence, the name no longer requires the ratification of Fries' authority. Early alternate names include Boletus solidus by English naturalist James Sowerby in 1809, and Gray's Leccinum edule. Gray's transfer of the species to Leccinum was later determined to be inconsistent with the rules of botanical nomenclature, and he apparently was unfamiliar with the earlier works of Fries when he published his arrangement of bolete species.
Boletus edulis is the type species of the genus Boletus. In Rolf Singer's classification of the Agaricales mushrooms, it is also the type species of section Boletus, a grouping of about 30 related boletes united by several characteristics: a mild-tasting, white flesh that does not change colour when exposed to air; a smooth to distinctly raised, netted pattern over at least the uppermost portion of the stem; a yellow-brown or olive-brown spore print; white tubes that later become yellowish then greenish, which initially appear to be stuffed with cotton; and cystidia that are not strongly coloured. Molecular analysis published in 1997 established that the bolete mushrooms are all derived from a common ancestor, and established the Boletales as an order separate from the Agaricales.
The generic name is derived from the Latin term bōlētus "mushroom", which was borrowed in turn from the Ancient Greek βωλίτης, "terrestrial fungus". Ultimately, this last word derives from bōlos/βῶλος "lump", "clod", and, metaphorically, "mushroom". The βωλίτης of Galen, like the boletus of Latin writers like Martial, Seneca and Petronius, is often identified as the much prized Amanita caesarea. The specific epithet edulis in Latin means "eatable" or "edible".
Common names
Common names for B. edulis vary by region. The standard Italian name, porcino (pl. porcini), means porcine; fungo porcino, in Italian, echoes the term suilli, literally "hog mushrooms", a term used by the Ancient Romans and still in use in southern Italian terms for this species. The derivation has been ascribed to the resemblance of young fruit bodies to piglets, or to the fondness pigs have for eating them. It is also known as "king bolete". The English penny bun refers to its rounded brownish shape. The German name Steinpilz (stone mushroom) refers to the species' firm flesh. In Austria, it is called Herrenpilz, the "noble mushroom", while in Mexico, the Spanish name is panza, meaning "belly". Another Spanish name, rodellon, means "small round boulder", while the Dutch name eekhoorntjesbrood means "squirrel's bread". Russian names are belyy grib (:ru:белый гриб; "white mushroom" as opposed to less valuable "black mushrooms") and borovik (:ru:боровик; from bor—"pine forest").
The vernacular name cep is derived from the Catalan cep or its French name cèpe, although the latter is a generic term applying to several related species. In France, it is more fully cèpe de Bordeaux, derived from the Gascon cep "trunk" for its fat stalk, ultimately from the Latin cippus "stake". Ceppatello, ceppatello buono, ceppatello bianco, giallo leonato, ghezzo, and moreccio are names from Italian dialects, and ciurenys or surenys is another term in Catalan. The French-born King Charles XIV John popularised B. edulis in Sweden after 1818, and is honoured in the local vernacular name Karljohanssvamp, as well as the Danish name Karl Johan svamp. The monarch cultivated the fungus about his residence, Rosersberg Palace. The Finnish name is herkkutatti, from herkku 'delicacy', and tatti, 'bolete'.
Description
The cap of this mushroom is broad at maturity. Slightly sticky to touch, it is convex in shape when young and flattens with age. The colour is generally reddish-brown fading to white in areas near the margin, and continues to darken as it matures. The stipe, or stem, is in height, and up to thick—rather large in comparison to the cap; it is club-shaped, or bulges out in the middle. It is finely reticulate on the upper portion, but smooth or irregularly ridged on the lower part. The under surface of the cap is made of thin tubes, the site of spore production; they are deep, and whitish in colour when young, but mature to a greenish-yellow. The angular pores, which do not stain when bruised, are small—roughly 2 to 3 pores per millimetre. In youth, the pores are white and appear as if stuffed with cotton (which are actually mycelia); as they age, they change colour to yellow and later to brown. The spore print is olive brown. The flesh of the fruit body is white, thick and firm when young, but becomes somewhat spongy with age. When bruised or cut, it either does not change colour, or turns a very light brown or light red. Fully mature specimens can weigh about ; a huge specimen collected on the Isle of Skye, Scotland, in 1995 bore a cap of , with a stipe in height and wide, and weighed . A similarly sized specimen found in Poland in 2013 made international news.
Boletus edulis is considered one of the safest wild mushrooms to pick for the table, as few poisonous species closely resemble it, and those that do may be easily distinguished by careful examination. The most similar poisonous mushroom may be the devil's bolete (Rubroboletus satanas), which has a similar shape, but has a red stem and stains blue on bruising. It is often confused with the very bitter and unpalatable Tylopilus felleus, but can be distinguished by the reticulation on the stalk; in porcini, it is a whitish, net-like pattern on a brownish stalk, whereas it is a dark pattern on white in the latter. Porcini have whitish pores while the other has pink. If in doubt, tasting a tiny bit of flesh will yield a bitter taste. It can also resemble the "bolete-like" Gyroporus castaneus, which is generally smaller, and has a browner stem. Boletus huronensis, an uncommon mushroom of northeastern North America, is another recognized look-alike known to cause severe gastrointestinal disorders.
The spores are elliptical to spindle-shaped, with dimensions of 12–17 by 5–7 μm. The basidia, the spore-bearing cells, are produced in a layer lining the tubes, and arrange themselves so their ends are facing the center of the tube; this layer of cells is known technically as a hymenium. The basidia are thin-walled, mostly attached to four spores, and measure 25–30 by 8–10 μm. Another cell type present in the hymenium is the cystidia, larger sterile cells that protrude beyond the basidia into the lumen of the hymenium, and act as air traps, regulating humidity. B. edulis has pleurocystidia (cystidia located on the face of a pore) that are thin-walled, roughly spindle-shaped to ventricose, and measure 30–45 by 7–10 μm; the "stuffed" feature of the hymenium is caused by cheilocystidia—cells found on the edges of the pores. The hyphae of B. edulis do not have clamp connections.
Related species
Several similar brownish-coloured species are sometimes considered subspecies or forms of this mushroom. In Europe, in addition to B. edulis (or cèpe de Bordeaux), the most popular are:
Cèpe bronzé ("dark cep"; Boletus aereus), much rarer than B. edulis, is more highly regarded by gourmets, and consequently more expensive. Usually smaller than B. edulis, it is also distinctively darker in colour. It is especially suited to drying.
Cèpe des pins ("pine tree cep"; Boletus pinophilus or Boletus pinicola) grows among pine trees. Rarer than B. edulis, it is less appreciated by gourmets than the two other kinds of porcini, but remains a mushroom rated above most others.
Cèpe d'été ("summer cep"; Boletus reticulatus), also less common and found earlier.
Molecular phylogenetic analyses have proven these three are all distinctive and separate species; other taxa formerly believed to be unique species or subspecies, such as B. betulicola, B. chippewaensis, B. persoonii, B. quercicola and B. venturii, are now known to be part of a B. edulis species complex with a wide morphological, ecological and geographic range, and that the genetic variability in this complex is low. Similar molecular technology has been developed to rapidly and accurately identify B. edulis and other commercially important fungi.
Three divergent lineages found in Yunnan province in China that are commonly marketed and sold as B. edulis (and are actually more closely related to B. aereus) were described in 2013 as B. bainiugan, B. meiweiniuganjun and B. shiyong. The classification has since been updated and expanded. All lineages are still members of Boletus sect. Boletus, the sensu sticto "porcini clade" of the genus.
Western North America has several species closely related to B. edulis. The white king bolete (Boletus barrowsii), found in parts of Colorado, New Mexico, Arizona, and California (and possibly elsewhere), is named after its discoverer Chuck Barrows. It is lighter in colour than B. edulis, having a cream-coloured cap with pink tones; often mycorrhizal with Ponderosa pine, it tends to grow in areas where there is less rainfall. Some find its flavour as good as if not better than B. edulis. The California king bolete (Boletus edulis var. grandedulis) can reach massive proportions, and is distinguished from B. edulis by a mature pore surface that is brown to slightly reddish. The cap colour appears to be affected by the amount of light received during its development, and may range from white in young specimens grown under thick canopy, to dark-brown, red-brown or yellow brown in those specimens receiving more light. The queen bolete (Boletus regineus), formerly considered a variety of B. aereus, is also a choice edible. It is generally smaller than B. edulis, and unlike that species, is typically found in mixed forests. The spring king bolete (Boletus rex-veris), formerly considered a variety of B. edulis or B. pinophilus, is found throughout western North America. In contrast to B. edulis, B. rex-veris tends to fruit in clusters, and, as its common name suggests, appears in the spring. B. fibrillosus is edible but considered inferior in taste.
Habitat and distribution
The fruit bodies of Boletus edulis can grow singly or in small clusters of two or three specimens. The mushroom's habitat consists of areas dominated by pine (Pinus spp.), spruce (Picea spp.), hemlock (Tsuga spp.) and fir (Abies spp.) trees, although other hosts include chestnut, chinquapin, beech, Keteleeria spp., Lithocarpus spp., and oak. In California, porcini have been collected in a variety of forests, such as coastal forests, dry interior oak forests and savannas and interior high-elevation montane mixed forests, to an altitude of . In northwestern Spain, they are common in scrublands dominated by the rock rose species Cistus ladanifer and Halimium lasianthum. In the Midi region of south-west France, they are especially favoured and locally called cèpe de Bordeaux after the town from which they are traded to the north and abroad.
Boletus edulis has a cosmopolitan distribution, concentrated in cool-temperate to subtropical regions. It is common in Europe—from northern Scandinavia, south to the extremities of Greece and Italy—and North America, where its southern range extends as far south as Mexico. It is well known from the Borgotaro area of Parma, Italy, and has PGI status there. The European distribution extends north to Scandinavia and south to southern Italy and Morocco. In the American Pacific Northwest, it can be found from May to October. In China, the mushroom can be found from the northeastern Heilongjiang to the Yunnan–Guizhou Plateau and Tibet. It has been recorded growing under Pinus and Tsuga in Sagarmatha National Park in Nepal, as well as in the Indian forests of Arunachal Pradesh. In West Asia, the species has been reported from the northwest forests of Iran.
Cultivation
Some steps have been made towards cultivating Boletus edulis, including mycorrhization of rockrose shrubs enhanced by helper bacteria.
Non-native introductions
Boletus edulis grows in some areas where it is not believed to be indigenous. It is often found underneath oak and silver birch in Hagley Park in central Christchurch, New Zealand, where it is likely to have been introduced, probably on the roots of container-grown beech, birch, and oak in the mid-19th century—around the time exotic trees began to be planted in the Christchurch area. Similarly, it has been collected in Adelaide Hills region of Australia in association with three species of introduced trees. It has been growing plentifully in association with pine forests in the southern KwaZulu-Natal Midlands in South Africa for more than 50 years and is believed to have been introduced with the import of pine trees. It also grows in pine plantations in neighboring Zimbabwe.
Ecology
Fruit body production
Italian folklore holds that porcini sprout up at the time of the new moon; research studies have tried to investigate more scientifically the factors that influence the production of fruit bodies. Although fruit bodies may appear any time from summer to autumn (June to November in the UK), their growth is known to be triggered by rainfall during warm periods of weather followed by frequent autumn rain with a drop in soil temperature. Above average rainfall may result in the rapid appearance of large numbers of boletes, in what is known in some circles as a "bolete year". A 2004 field study indicated that fruit body production is enhanced by an open and sunny wood habitat, corroborating an earlier observation made in a Zimbabwean study; removal of the litter layer on the forest floor appeared to have a negative effect on fruit body production, but previous studies reported contradictory results. A Lithuanian study conducted in 2001 concluded that the maximal daily growth rate of the cap (about 21 mm or 0.8 in) occurred when the relative air humidity was the greatest, and the fruit bodies ceased growing when the air humidity dropped below 40%. Factors most likely to inhibit the appearance of fruit bodies included prolonged drought, inadequate air and soil humidity, sudden decreases of night air temperatures, and the appearance of the first frost. Plots facing north tend to produce more mushrooms compared to equivalent plots facing south.
Mycorrhizal associations
Boletus edulis is mycorrhizal—it is in a mutualistic relationship with the roots of plants (hosts), in which the fungus exchanges nitrogen and other nutrients extracted from the environment for fixed carbon from the host. Other benefits for the plant are evident: in the case of the Chinese chestnut, the formation of mycorrhizae with B. edulis increases the ability of plant seedlings to resist water stress, and increases leaf succulence, leaf area, and water-holding ability. The fungus forms a sheath of tissue around terminal, nutrient-absorbing root tips, often inducing a high degree of branching in the tips of the host, and penetrating into the root tissue, forming, to some mycologists, the defining feature of ectomycorrhizal relationships, a hartig net. The ectomycorrhizal fungi are then able to exchange nutrients with the plant, effectively expanding the root system of the host plant to the furthest reaches of the symbiont fungi. Compatible hosts may belong to multiple families of vascular plants that are widely distributed throughout the Northern Hemisphere; according to one 1995 estimate, there are at least 30 host plant species distributed over more than 15 genera. Examples of mycorrhizal associates include Chinese red pine, Mexican weeping pine, Scots pine, Norway spruce, Coast Douglas-fir, mountain pine, and Virginia pine. The fungus has also been shown to associate with gum rockrose, a pioneer early stage shrub that is adapted for growth in degraded areas, such as burned forests. These and other rockrose species are ecologically important as fungal reservoirs, maintaining an inoculum of mycorrhizal fungi for trees that appear later in the forest regrowth cycle.
The mushroom has been noted to often co-occur with Amanita muscaria or A. rubescens, although it is unclear whether this is due to a biological association between the species, or because of similarities in growing season, habitat, and ecological requirements. An association has also been reported between B. edulis and Amanita excelsa on Pinus radiata ectomycorrhizae in New Zealand, suggesting that other fungi may influence the life cycle of porcini. A 2007 field study revealed little correlation between the abundance of fruit bodies and presence of its mycelia below ground, even when soil samples were taken from directly beneath the mushroom; the study concluded that the triggers leading to formation of mycorrhizae and production of the fruit bodies were more complex.
Heavy-metal contamination
Boletus edulis is known to be able to tolerate and even thrive on soil that is contaminated with toxic heavy metals, such as soil that might be found near metal smelters. The mushroom's resistance to heavy-metal toxicity is conferred by a biochemical called a phytochelatin—an oligopeptide whose production is induced after exposure to metal. Phytochelatins are chelating agents, capable of forming multiple bonds with the metal; in this state, the metal cannot normally react with other elements or ions and is stored in a detoxified form in the mushroom tissue.
Pests and predators
The fruit bodies of B. edulis can be infected by the parasitic mould-like fungus Hypomyces chrysospermus, known as the bolete eater, which manifests itself as a white, yellow, or reddish-brown cottony layer over the surface of the mushroom. Some reported cases of stomach ache following consumption of dried porcini have been attributed to the presence of this mould on the fruit bodies. The mushroom is also used as a food source by several species of mushroom flies, as well as other insects and their larvae. An unidentified species of virus was reported to have infected specimens found in the Netherlands and in Italy; fruit bodies affected by the virus had relatively thick stems and small or no caps, leading to the name "little-cap disease".
Boletus edulis is a food source for animals such as the banana slug (Ariolimax columbianus), the long-haired grass mouse, the red squirrel, and, as noted in one isolated report, the fox sparrow.
Culinary uses
Boletus edulis, as the species epithet edulis () indicates, is an edible mushroom. Italian chef and restaurateur Antonio Carluccio has described it as representing "the wild mushroom par excellence", and hails it as the most rewarding of all fungi in the kitchen for its taste and versatility. Considered a choice edible, particularly in France, Germany, Poland and Italy, it was widely written about by the Roman writers Pliny the Elder and Martial, although ranked below the esteemed Amanita caesarea. When he was served suilli instead of boleti, the disgruntled Martial wrote:
The flavour of porcini has been described as nutty and slightly meaty, with a smooth, creamy texture, and a distinctive aroma reminiscent of sourdough. Young, small porcini are most appreciated, as the large ones often harbour maggots (insect larvae), and become slimy, soft and less tasty with age. The fruit bodies are collected by holding the stipe near the base and twisting gently. Cutting the stipe with a knife may risk the part left behind rotting and the mycelium being destroyed. Peeling and washing are not recommended. The fruit bodies are highly perishable, due largely to the high water content (around 90%), the high level of enzyme activity, and the presence of a flora of microorganisms. Caution should be exercised when collecting specimens from potentially polluted or contaminated sites, as several studies have shown that the fruit bodies can bioaccumulate toxic heavy metals like mercury, cadmium, caesium and polonium. Bioaccumulated metals or radioactive fission decay products are like chemical signatures: chemical and radiochemical analysis can be used to identify the origin of imported specimens, and for long-term radioecological monitoring of polluted areas.
Porcini are sold fresh in markets in summer and autumn in Europe and Russia, and dried or canned at other times of the year, and distributed worldwide to countries where they are not otherwise found. They are eaten and enjoyed raw, sautéed with butter, ground into pasta, in soups, and in many other dishes. In France, they are used in recipes such as cèpes à la Bordelaise, cèpe frits and cèpe aux tomates. Porcini risotto is a traditional Italian autumn dish. Porcini are a feature of many cuisines, including Provençal, and Viennese. In Thailand they are used in soups and consumed blanched in salads. Porcini can also be frozen, either while raw or after cooking in butter. The colour, aroma, and taste of porcini deteriorate noticeably after being frozen for four months. Blanching (or soaking and blanching) as a processing step before freezing can extend the freezer life to 12 months. They are also one of the few species sold commercially as pickled mushrooms.
Dried
Boletus edulis is well suited to drying—its flavour intensifies, it is easily reconstituted, and its resulting texture is pleasant. Reconstitution is done by soaking in hot, but not boiling, water for about twenty minutes; the water used is infused with the mushroom aroma and it too can be used in subsequent cooking. Dried porcini have more protein than most other commonly consumed vegetables, apart from soybeans. Some of their protein is indigestible, though digestibility is improved with cooking.
Like other boletes, porcini can be dried by being strung separately on twine and hung close to the ceiling of a kitchen. Alternatively, the mushrooms can be dried by cleaning with a brush (washing is not recommended), and then placing them in a wicker basket or bamboo steamer on top of a boiler or hot water tank. Another method is drying in an oven at for two to three hours, then increasing the temperature to until crisp or brittle. Once dry, they must be kept in an airtight container. Importantly for commercial production, porcini retain their flavour after industrial preparation in a pressure cooker or after canning or bottling, and are thus useful for manufacturers of soups or stews. The addition of a few pieces of dried porcino can significantly add to flavour, and they are a major ingredient of the pasta sauce known as carrettiera (carter's sauce). The drying process is known to induce the formation of various volatile substances that contribute to the mushroom's aroma. Chemical analysis has shown that the odour of the dried mushroom is a complex mixture of 53 volatile compounds.
Commercial harvest
A 1998 estimate suggested that the total annual worldwide consumption of Boletus edulis and closely related species (B. aereus, B. pinophilus, and B. reticulatus) was between 20,000 and 100,000 tons. Approximately 2,700 tonnes (3,000 tons) were sold in France, Italy and Germany in 1988, according to official figures. The true amount consumed far exceeds this, as the official sales figures did not account for informal sales or consumption by collectors. They are widely exported and sold in dried form, reaching countries where they do not occur naturally, such as Australia and New Zealand. The autonomous community of Castile and León in Spain produces 7,700 tonnes (8,500 tons) annually. In autumn, the price of porcini in the Northern Hemisphere typically ranges between $20 and $80 per kilogram, although in New York in 1997 the wholesale price rose to more than $200 per kilogram due to scarcity.
In the vicinity of Borgotaro in the Province of Parma of northern Italy, the four species Boletus edulis, B. aereus, B. aestivalis and B. pinophilus have been recognised for their superior taste and officially termed Fungo di Borgotaro. Here these mushrooms have been collected for centuries and exported commercially. Owing to the globalisation of the mushroom trade most of the porcini commercially available in Italy or exported by Italy no longer originate there. Porcini and other mushrooms are also imported into Italy from various locations, especially China and eastern European countries; these are then often re-exported under the "Italian porcini" label.
In Italy the disconnect with local production has had an adverse effect on quality; for example in the 1990s some of the dried porcino mushrooms exported to Italy from China contained species of genus Tylopilus, which are rather similar in appearance and when dried are difficult for both mushroom labourers and mycologists alike to distinguish from Boletus. Tylopilus species typically have a very bitter taste, which is imparted to the flavour of the porcini with which they are mixed.
After the fall of the Iron Curtain and the subsequent reduction of economic and political barriers, central and eastern European countries with local mushroom harvesting traditions, such as Albania, Bulgaria, Macedonia, Romania, Serbia and Slovenia, developed into exporters of porcini, concentrating primarily on the Italian market. Porcini and other wild fungi from these countries are also destined for France, Germany and other western European markets, where demand for them exists but collection on a commercial scale does not. Picking B. edulis has become an annual seasonal income earner and pastime in countries like Bulgaria, especially for many Roma communities and the unemployed. A lack of control of the harvest has led to heavy exploitation of the mushroom resource.
Like many other strictly mycorrhizal fungi, B. edulis has eluded cultivation attempts for years. The results of some studies suggest that unknown components of the soil microflora might be required for B. edulis to establish a mycorrhizal relationship with the host plant. Successful attempts at cultivating B. edulis have been made by Spanish scientists by mycorrhization of Cistus species, with Pseudomonas fluorescens bacteria helping the mycorrhiza.
Nutrition
Boletus edulis mushrooms are 9% carbohydrates, 3% fat, and 7% protein (table). Fresh mushrooms consist of over 80% moisture, although reported values tend to differ somewhat as moisture content can be affected by environmental temperature and relative humidity during growth and storage. The carbohydrate component contains the monosaccharides glucose, mannitol and α,α-trehalose, the polysaccharide glycogen, and the water-insoluble structural polysaccharide chitin, which accounts for up to 80–90% of dry matter in mushroom cell walls. Chitin, hemicellulose, and pectin-like carbohydrates—all indigestible by humans—contribute to the high proportion of insoluble fibre in B. edulis.
The total lipid, or crude fat, content makes up 3% of the dry matter of the mushroom. The proportion of fatty acids (expressed as a % of total fatty acids) are: linoleic acid 42%, oleic acid 36%, palmitic acid 10%, and stearic acid 3%.
A comparative study of the amino acid composition of eleven Portuguese wild edible mushroom species showed Boletus edulis to have the highest total amino acid content.
B. edulis mushrooms are rich in the dietary minerals, sodium, iron, calcium, and magnesium, with amounts varying according to the mushroom component and to soil composition in the geographic region of China where they were sampled. They also have high content of B vitamins and tocopherols. B. edulis contains appreciable amounts of selenium, a trace mineral, although the bioavailability of mushroom-derived selenium is low.
Phytochemicals and research
Boletus edulis fruit bodies contain diverse phytochemicals, including 500 mg of ergosterol per 100 g of dried mushroom, and ergothioneine. The fruit bodies contain numerous polyphenols, especially a high content of rosmarinic acid, and organic acids (such as oxalic, citric, malic, succinic and fumaric acids), and alkaloids.
Aroma
Aroma compounds giving B. edulis mushrooms their characteristic fragrance include some 100 components, such as esters and fatty acids. In a study of aroma compounds, 1-octen-3-one was the most prevalent chemical detected in raw mushrooms, with pyrazines having increased aroma effect and elevated content after drying.
| Biology and health sciences | Edible fungi | null |
428052 | https://en.wikipedia.org/wiki/Sagittarius%20Dwarf%20Spheroidal%20Galaxy | Sagittarius Dwarf Spheroidal Galaxy | The Sagittarius Dwarf Spheroidal Galaxy (Sgr dSph), also known as the Sagittarius Dwarf Elliptical Galaxy (Sgr dE or Sag DEG), is an elliptical loop-shaped satellite galaxy of the Milky Way. It contains four globular clusters in its main body, with the brightest of them — NGC 6715 (M54) — known well before the discovery of the galaxy itself in 1994. Sgr dSph is roughly 10,000 light-years in diameter, and is currently about 70,000 light-years from Earth, travelling in a polar orbit (an orbit passing over the Milky Way's galactic poles) at a distance of about 50,000 light-years from the core of the Milky Way (about one third of the distance of the Large Magellanic Cloud). In its looping, spiraling path, it has passed through the plane of the Milky Way several times in the past. In 2018, the Gaia project of the European Space Agency showed that Sgr dSph had caused perturbations in a set of stars near the Milky Way's core, causing unexpected rippling movements of the stars triggered when it moved through the Milky Way between 300 and 900 million years ago.
Features
Officially discovered in 1994, by Rodrigo Ibata, Mike Irwin, and Gerry Gilmore, Sgr dSph was immediately recognized as being the nearest known neighbor to the Milky Way at the time. (The disputed Canis Major Dwarf Galaxy, discovered in 2003, might be the actual nearest neighbor.) Although it is one of the closest companion galaxies to the Milky Way, the main parent cluster is on the opposite side of the Galactic Center from Earth, and consequently is very faint, although covering a large area of the sky. Sgr dSph appears to be an older galaxy with little interstellar dust, composed largely of Population II stars, older and metal-poor, as compared to the Milky Way. No neutral hydrogen gas related to Sgr dSph has been found.
Further discoveries by astrophysics teams from both the University of Virginia and the University of Massachusetts Amherst, drawing upon the 2MASS Two-Micron All Sky Infrared Survey data, revealed the entire loop-shaped structure. In 2003 with the aid of infrared telescopes and super computers, Steven Majewski, Michael Skrutskie, and Martin Weinberg were able to help create a new star map, picking out the full Sagittarius Dwarf presence, position, and looping shape from the mass of background stars and finding this smaller galaxy to be at a near right angle to the plane of the Milky Way.
Globular clusters
Sgr dSph has at least nine known globular clusters. One, M 54, appears to reside at its core, while three others reside within the main body of the galaxy: Terzan 7, Terzan 8 and Arp 2.
Additionally, Palomar 12, Whiting 1, NGC 2419, NGC 4147, and NGC 5634 are found within its extended stellar streams. However, this is an unusually low number of globular clusters, and an analysis of VVV and Gaia EDR3 data has found at least twenty more. The newly discovered globular clusters tend to be more metal-rich than previously known globular clusters.
Metallicity
Sgr dSph has multiple stellar populations, ranging in age from the oldest globular clusters (almost as old as the universe itself) to trace populations as young as several hundred million years (mya). It also exhibits an age-metallicity relationship, in that its old populations are metal poor () while its youngest populations have super-solar abundances.
Geometry and dynamics
Based on its current trajectory, the Sgr dSph main cluster is about to pass through the galactic disc of the Milky Way within the next hundred million years, while the extended loop-shaped ellipse is already extended around and through our local space and on through the Milky Way galactic disc, and in the process of slowly being absorbed into the larger galaxy, calculated at 10,000 times the mass of Sgr dSph. The dissipation of the Sgr dSph main cluster and its merger with the Milky Way stream is expected to be complete within a billion years from now.
At first, many astronomers thought that Sgr dSph had already reached an advanced state of destruction, so that a large part of its original matter was already mixed with that of the Milky Way. However, Sgr dSph still has coherence as a dispersed elongated ellipse, and appears to move in a roughly polar orbit around the Milky Way as close as 50,000 light-years from the galactic core. Although it may have begun as a spherical object before falling towards the Milky Way, Sgr dSph is now being torn apart by immense tidal forces over hundreds of millions of years. Numerical simulations suggest that stars ripped out from the dwarf would be spread out in a long stellar stream along its path, which were subsequently detected.
However, some astronomers contend that Sgr dSph has been in orbit around the Milky Way for some billions of years, and has already orbited it approximately ten times. Its ability to retain some coherence despite such strains would indicate an unusually high concentration of dark matter within that galaxy.
In 1999, Johnston et al. concluded that Sgr dSph has orbited the Milky Way for at least one gigayear and that during that time its mass has decreased by a factor of two or three. Its orbit is found to have galactocentric distances that oscillate between ≈13 and ≈41 kpc with a period of 550 to 750 million years. The last perigalacticon was approximately fifty million years ago. Also in 1999, Jiang & Binney found that it may have started its infall into the Milky Way at a point more than 200 kpc away if its starting mass was as large as ≈1011.
The models of both its orbit and the Milky Way's potential field could be improved by proper motion observations of Sgr dSph's stellar debris. This issue is under intense investigation, with computational support by the MilkyWay@Home project.
A simulation published in 2011 suggested that the Milky Way may have obtained its spiral structure as a result of repeated collisions with Sgr dSph.
In 2018, the Gaia project of the European Space Agency, designed primarily to investigate the origin, evolution and structure of the Milky Way, delivered the largest and most precise census of positions, velocities and other stellar properties of more than a billion stars, which showed that Sgr dSph had caused perturbations in a set of stars near the Milky Way's core, causing unexpected rippling movements of the stars triggered when it sailed past the Milky Way between 300 and 900 million years ago.
A 2019 study by TCU Graduate Student Matthew Melendez and co-authors concluded that Sgr dSph had a decreasing metallicity trend as a function of radius, with a larger spread in metallicity in the core relative to the outer regions. Also, they did find evidence for the first time for two distinct populations in alpha abundances as a function of metallicity.
A 2020 study concluded that collisions between the Sagittarius Dwarf Spheroidal Galaxy and the Milky Way triggered major episodes of star formation in the latter, based on data taken from the Gaia project.
| Physical sciences | Notable galaxies | Astronomy |
428075 | https://en.wikipedia.org/wiki/Orb-weaver%20spider | Orb-weaver spider | Orb-weaver spiders are members of the spider family Araneidae. They are the most common group of builders of spiral wheel-shaped webs often found in gardens, fields, and forests. The English word "orb" can mean "circular", hence the English name of the group. Araneids have eight similar eyes, hairy or spiny legs, and no stridulating organs.
The family has a cosmopolitan distribution, including many well-known large or brightly colored garden spiders. With 3,108 species in 186 genera worldwide, the Araneidae comprise one of the largest family of spiders (with the Salticidae and Linyphiidae). Araneid webs are constructed in a stereotypical fashion, where a framework of nonsticky silk is built up before the spider adds a final spiral of silk covered in sticky droplets.
Orb webs are also produced by members of other spider families. The long-jawed orb weavers (Tetragnathidae) were formerly included in the Araneidae; they are closely related, being part of the superfamily Araneoidea. The family Arkyidae has been split off from the Araneidae. The cribellate or hackled orb-weavers (Uloboridae) belong to a different group of spiders. Their webs are strikingly similar, but use a different kind of silk.
Description
Generally, orb-weaving spiders are three-clawed builders of flat webs with sticky spiral capture silk. The building of a web is an engineering feat, begun when the spider floats a line on the wind to another surface. The spider secures the line and then drops another line from the center, making a "Y". The rest of the scaffolding follows with many radii of nonsticky silk being constructed before a final spiral of sticky capture silk.
The third claw is used to walk on the nonsticky part of the web. Characteristically, the prey insect that blunders into the sticky lines is stunned by a quick bite, and then wrapped in silk. If the prey is a venomous insect, such as a wasp, wrapping may precede biting and/or stinging. Much of the orb-spinning spiders' success in capturing insects depends on the web not being visible to the prey, with the stickiness of the web increasing the visibility, thus decreasing the chances of capturing prey. This leads to a trade-off between the visibility of the web and the web's prey-retention ability.
Many orb-weavers build a new web each day. Most orb-weavers tend to be active during the evening hours; they hide for most of the day. Generally, towards evening, the spider consumes the old web, rests for about an hour, then spins a new web in the same general location. Thus, the webs of orb-weavers are generally free of the accumulation of detritus common to other species, such as black widow spiders.
Some orb-weavers do not build webs at all. Members of the genera Mastophora in the Americas, Cladomelea in Africa, and Ordgarius in Australia produce sticky globules, which contain a pheromone analog. The globule is hung from a silken thread dangled by the spider from its front legs. The pheromone analog attracts male moths of only a few species. These get stuck on the globule and are reeled in to be eaten. Both genera of bolas spiders are highly camouflaged and difficult to locate.
In the Araneus diadematus, variables such as wind, web support, temperatures, humidity, and silk supply all proved to be variables in web construction. When studied against the tests of nature, the spiders were able to decide what shape to make their web, how many capture spirals, or the width of their web. Though it could be expected for these spiders to just know these things, it is not well researched yet as to just how the arachnid knows how to change their web design based on their surroundings. Some scientists suggest that it could be through the spider's spatial learning on their environmental surroundings and the knowing of what will or will not work compared to natural behavioristic rules.
The spiny orb-weaving spiders in the genera Gasteracantha and Micrathena look like plant seeds or thorns hanging in their orb-webs. Some species of Gasteracantha have very long, horn-like spines protruding from their abdomens.
One feature of the webs of some orb-weavers is the stabilimentum, a crisscross band of silk through the center of the web. It is found in several genera, but Argiope – the yellow and banded garden spiders of North America – is a prime example. As orb-weavers age, they tend to have less production of their silk; many adult orb-weavers can then depend on their coloration to attract more of their prey. The band may be a lure for prey, a marker to warn birds away from the web, and a camouflage for the spider when it sits in the web. The stabilimentum may decrease the visibility of the silk to insects, thus making it harder for prey to avoid the web. The orb-web consists of a frame and supporting radii overlaid with a sticky capture spiral, and the silks used by orb-weaver spiders have exceptional mechanical properties to withstand the impact of flying prey. The orb-weaving spider Zygiella x-notata produces a unique orb-web with a characteristic missing sector, similar to other species of the Zygiella genus in the Araneidae family.
During the Cretaceous, a radiation of flowering plants and their insect pollinators occurred. Fossil evidence shows that the orb web was in existence at this time, which permitted a concurrent radiation of the spider predators along with their insect prey. The capacity of orb–webs to absorb the impact of flying prey led orbicularian spiders to become the dominant predators of aerial insects in many ecosystems. Insects and spiders have comparable rates of diversification, suggesting they co-radiated, and the peak of this radiation occurred 100 Mya, before the origin of angiosperms. Vollrath and Selden (2007) make the bold proposition that insect evolution was driven less by flowering plants than by spider predation – particularly through orb webs – as a major selective force. On the other hand some analyses have yielded estimates as high as 265 Mya, with a large number (including Dimitrov et al 2016) intermediate between the two.
Most arachnid webs are vertical and the spiders usually hang with their heads downward. A few webs, such as those of orb-weavers in the genus Metepeira, have the orb hidden within a tangled space of web. Some Metepiera species are semisocial and live in communal webs. In Mexico, such communal webs have been cut out of trees or bushes and used for living fly paper. In 2009, workers at a Baltimore wastewater treatment plant called for help to deal with over 100 million orb-weaver spiders, living in a community that managed to spin a phenomenal web that covered some 4 acres of a building, with spider densities in some areas reaching 35,176 spiders per cubic meter.
Taxonomy
The oldest known true orb-weaver is Mesozygiella dunlopi, from the Lower Cretaceous. Several fossils provide direct evidence that the three major orb-weaving families, namely the Araneidae, Tetragnathidae, and Uloboridae, had evolved by this time, about 140 Mya. They probably originated during the Jurassic (). Based on new molecular evidence in silk genes, all three families are likely to have a common origin.
The two superfamilies, Deinopoidea and Araneoidea, have similar behavioral sequences and spinning apparatuses to produce architecturally similar webs. The latter weave true viscid silk with an aqueous glue property, and the former use dry fibrils and sticky silk. The Deinopoidea (including the Uloboridae), have a cribellum – a flat, complex spinning plate from which the cribellate silk is released.
They also have a calamistrum – an apparatus of bristles used to comb the cribellate silk from the cribellum. The Araneoidea, or the "ecribellate" spiders, do not have these two structures. The two groups of orb-weaving spiders are morphologically very distinct, yet much similarity exists between their web forms and web construction behaviors. The cribellates retained the ancestral character, yet the cribellum was lost in the escribellates. The lack of a functional cribellum in araneoids is most likely synapomorphic.
If the orb-weaver spiders are a monophyletic group, the fact that only some species in the group lost a feature adds to the controversy. The cribellates are split off as a separate taxon that retained the primitive feature, which makes the lineage paraphyletic and not synonymous with any real evolutionary lineage. The morphological and behavioral evidence surrounding orb webs led to the disagreement over a single or a dual origin. While early molecular analysis provided more support for a monophyletic origin, other evidence indicates that orb-weavers evolved earlier phylogenetically than previously thought, and were extinct at least three times during the Cretaceous.
Reproduction
Araneid species either mate at the central hub of the web, where the male slowly traverses the web, trying not to get eaten, and when reaching the hub, mounts the female; or the male constructs a mating thread inside or outside the web to attract the female via vibratory courtship, and if successful, mating occurs on the thread.
In the cannibalistic and polyandrous orb-web spider Argiope bruennichi, the much smaller males are attacked during their first copulation and are cannibalized in up to 80% of the cases. All surviving males die after their second copulation, a pattern observed in other Argiope species. Whether a male survives his first copulation depends on the duration of the genital contact; males that jump off early (before 5 seconds) have a chance of surviving, while males that copulate longer (greater than 10 seconds) invariably die. Prolonged copulation, although associated with cannibalism, enhances sperm transfer and relative paternity.
When males mated with a nonsibling female, the duration of their copulation was prolonged, and consequently the males were cannibalized more frequently. When males mated with a sibling female, they copulated briefly, thus were more likely to escape cannibalism. By escaping, their chance of mating again with an unrelated female likely would be increased. These observations suggest that males can adaptively adjust their investment based on the degree of genetic relatedness of the female to avoid inbreeding depression.
Sexual size dimorphism
Sexual dimorphism refers to physical differences between males and females of the same species. One such difference can be in size.
Araneids often exhibit size dimorphism typically known as extreme sexual size dimorphism, due to the extent of differences in size. The size difference among species of Araneidae ranges greatly. Some females, such as those of the Nephila pilipes, can be at least 9 times larger than the male, while others are only slightly larger than the male. The larger size female is typically thought to be selected through fecundity selection, the idea that bigger females can produce more eggs, thus more offspring. Although a great deal of evidence points towards the greatest selection pressure on larger female size, some evidence indicates that selection can favor small male size, as well.
Araneids also exhibit a phenomenon called sexual cannibalism, which is commonly found throughout the Araneidae. Evidence suggests a negative correlation between sexual size dimorphism and instances of sexual cannibalism. Other evidence, however, has shown that differences in cannibalistic events among araneids when having smaller or slightly larger males is advantageous.
Some evidence has shown that extreme dimorphism may be the result of males avoiding detection by the females. For males of these species, being smaller in size may be advantageous in moving to the central hub of a web so female spiders may be less likely to detect the male, or even if detected as prey to be eaten, the small size may indicate little nutritional value. Larger-bodied male araneids may be advantageous when mating on a mating thread because the thread is constructed from the edge of the web orb to structural threads or to nearby vegetation. Here larger males may be less likely to be cannibalized, as the males are able to copulate while the female is hanging, which may make them safer from cannibalism. In one subfamily of Araneid that uses a mating thread, Gasteracanthinae, sexual cannibalism is apparently absent despite extreme size dimorphism.
Genera
, the World Spider Catalog accepts the following genera:
Abba Castanheira & Framenau, 2023 – Australia (Queensland, New South Wales)
Acacesia Simon, 1895 — South America, North America
Acantharachne Tullgren, 1910 — Congo, Madagascar, Cameroon
Acanthepeira Marx, 1883 — North America, Brazil, Cuba
Acroaspis Karsch, 1878 — New Zealand, Australia
Acrosomoides Simon, 1887 — Madagascar, Cameroon, Congo
Actinacantha Simon, 1864 — Indonesia
Actinosoma Holmberg, 1883 — Colombia, Argentina
Aculepeira Chamberlin & Ivie, 1942 — North America, Central America, South America, Asia, Europe
Acusilas Simon, 1895 — Asia
Aethriscus Pocock, 1902 — Congo
Aethrodiscus Strand, 1913 — Central Africa
Aetrocantha Karsch, 1879 — Central Africa
Afracantha Dahl, 1914 — Africa
Agalenatea Archer, 1951 — Ethiopia, Asia
Alenatea Song & Zhu, 1999 — Asia
Allocyclosa Levi, 1999 — United States, Panama, Cuba
Alpaida O. Pickard-Cambridge, 1889 — Central America, South America, Mexico, Caribbean
Amazonepeira Levi, 1989 — South America
Anepsion Strand, 1929 — Oceania, Asia
Aoaraneus Tanikawa, Yamasaki & Petcharad, 2021 — China, Japan, Korea, Taiwan
Arachnura Clerck, 1863
Araneus Clerck, 1757
Araniella Chamberlin & Ivie, 1942 — Asia
Aranoethra Butler, 1873 — Africa
Argiope Audouin, 1826 — Asia, Oceania, Africa, North America, South America, Costa Rica, Cuba, Portugal
Artifex Kallal & Hormiga, 2018 — Australia
Artonis Simon, 1895 — Myanmar, Ethiopia
Aspidolasius Simon, 1887 — South America
Augusta O. Pickard-Cambridge, 1877 — Madagascar
Austracantha Dahl, 1914 — Australia
Backobourkia Framenau, Dupérré, Blackledge & Vink, 2010 — Australia, New Zealand
Bertrana Keyserling, 1884 — South America, Central America
Bijoaraneus Tanikawa, Yamasaki & Petcharad, 2021 — Africa, Asia, Oceania
Caerostris Thorell, 1868 — Africa, Asia
Carepalxis L. Koch, 1872 — Oceania, South America, Mexico, Jamaica
Celaenia Thorell, 1868 — Australia, New Zealand
Cercidia Thorell, 1869 — Russia, Kazakhstan, India
Chorizopes O. Pickard-Cambridge, 1871 — Asia, Madagascar
Chorizopesoides Mi & Wang, 2018 — China, Vietnam
Cladomelea Simon, 1895 — South Africa, Congo
Clitaetra Simon, 1889 — Africa, Sri Lanka
Cnodalia Thorell, 1890 — Indonesia, Japan
Coelossia Simon, 1895 — Sierra Leone, Mauritius, Madagascar
Colaranea Court & Forster, 1988 — New Zealand
Collina Urquhart, 1891 — Australia
Colphepeira Archer, 1941 — United States, Mexico
Courtaraneus Framenau, Vink, McQuillan & Simpson, 2022 — New Zealand
Cryptaranea Court & Forster, 1988 — New Zealand
Cyclosa Menge, 1866 — Caribbean, Asia, Oceania, South America, North America, Central America, Africa, Europe
Cyphalonotus Simon, 1895 — Asia, Africa
Cyrtarachne Thorell, 1868 — Asia, Africa, Oceania
Cyrtobill Framenau & Scharff, 2009 — Australia
Cyrtophora Simon, 1864 — Asia, Oceania, Dominican Republic, Costa Rica, South America, Africa
Deione Thorell, 1898 — Myanmar
Deliochus Simon, 1894 — Australia, Papua New Guinea
Dolophones Walckenaer, 1837 — Australia, Indonesia
Dubiepeira Levi, 1991 — South America
Edricus O. Pickard-Cambridge, 1890 — Mexico, Panama, Ecuador
Enacrosoma Mello-Leitão, 1932 — South America, Central America, Mexico
Encyosaccus Simon, 1895 — South America
Epeiroides Keyserling, 1885 — Costa Rica, Brazil
Eriophora Simon, 1864 — Oceania, United States, South America, Central America, Africa
Eriovixia Archer, 1951 — Asia, Papua New Guinea, Africa
Eustacesia Caporiacco, 1954 — French Guiana
Eustala Simon, 1895 — South America, North America, Central America, Caribbean
Exechocentrus Simon, 1889 — Madagascar
Faradja Grasshoff, 1970 — Congo
Friula O. Pickard-Cambridge, 1897 — Indonesia
Galaporella Levi, 2009 — Ecuador
Gasteracantha Sundevall, 1833 — Oceania, Asia, United States, Africa, Chile
Gastroxya Benoit, 1962 — Africa
Gea C. L. Koch, 1843 — Africa, Oceania, Asia, United States, Argentina
Gibbaranea Archer, 1951 — Asia, Europe, Algeria
Glyptogona Simon, 1884 — Sri Lanka, Italy, Israel
Gnolus Simon, 1879 — Chile, Argentina
Guizygiella Zhu, Kim & Song, 1997 — Asia
Herennia Thorell, 1877 — Asia, Oceania
Heterognatha Nicolet, 1849 — Chile
Heurodes Keyserling, 1886 — Asia, Australia
Hingstepeira Levi, 1995 — South America
Hortophora Framenau & Castanheira, 2021 — Oceania
Hypognatha Guérin, 1839 — South America, Central America, Mexico, Trinidad
Hypsacantha Dahl, 1914 — Africa
Hypsosinga Ausserer, 1871 — Asia, North America, Greenland, Africa
Ideocaira Simon, 1903 — South Africa
Indoetra Kuntner, 2006 — Sri Lanka
Isoxya Simon, 1885 — Africa, Yemen
Kaira O. Pickard-Cambridge, 1889 — North America, South America, Cuba, Guatemala
Kangaraneus Castanheira & Framenau, 2023 — Australia
Kapogea Levi, 1997 — Mexico, South America, Central America
Kilima Grasshoff, 1970 — Congo, Seychelles, Yemen
Larinia Simon, 1874 — Asia, Africa, South America, Europe, Oceania, North America
Lariniaria Grasshoff, 1970 — Asia
Larinioides Caporiacco, 1934 — Asia
Lariniophora Framenau, 2011 — Australia
Leviana Framenau & Kuntner, 2022 — Australia
Leviaraneus Tanikawa & Petcharad, 2023 — Asia
Leviellus Wunderlich, 2004 — Asia, France
Lewisepeira Levi, 1993 — Panama, Mexico, Jamaica
Lipocrea Thorell, 1878 — Asia, Europe
Macracantha Simon, 1864 — India, China, Indonesia
Madacantha Emerit, 1970 — Madagascar
Mahembea Grasshoff, 1970 — Central and East Africa
Mangora O. Pickard-Cambridge, 1889 — Asia, North America, South America, Central America, Caribbean
Mangrovia Framenau & Castanheira, 2022 — Australia
Manogea Levi, 1997 — South America, Central America, Mexico
Mastophora Holmberg, 1876 — South America, North America, Central America, Cuba
Mecynogea Simon, 1903 — North America, South America, Cuba
Megaraneus Lawrence, 1968 — Africa
Melychiopharis Simon, 1895 — Brazil
Metazygia F. O. Pickard-Cambridge, 1904 — South America, Central America, North America, Caribbean
Metepeira F. O. Pickard-Cambridge, 1903 — North America, Caribbean, South America, Central America
Micrathena Sundevall, 1833 — South America, Caribbean, Central America, North America
Micrepeira Schenkel, 1953 — South America, Costa Rica
Micropoltys Kulczyński, 1911 — Papua New Guinea, Australia
Milonia Thorell, 1890 — Singapore, Indonesia, Myanmar
Molinaranea Mello-Leitão, 1940 — Chile, Argentina
Nemoscolus Simon, 1895 — Africa
Nemosinga Caporiacco, 1947 — Tanzania
Nemospiza Simon, 1903 — South Africa
Neogea Levi, 1983 — Papua New Guinea, India, Indonesia
Neoscona Simon, 1864 — Asia, Africa, Europe, Oceania, North America, Cuba, South America
Nephila Leach, 1815 — Asia, Oceania, United States, Africa, South America
Nephilengys L. Koch, 1872 — Asia, Oceania
Nephilingis Kuntner, 2013 — South America, Africa
Nicolepeira Levi, 2001 — Chile
Novakiella Court & Forster, 1993 — Australia, New Zealand
Novaranea Court & Forster, 1988 — Australia, New Zealand
Nuctenea Simon, 1864 — Algeria, Asia, Europe
Oarces Simon, 1879 — Brazil, Chile, Argentina
Ocrepeira Marx, 1883 — South America, Central America, Caribbean, North America
Ordgarius Keyserling, 1886 — Asia, Oceania
Paralarinia Grasshoff, 1970 — Congo, South Africa
Paraplectana Brito Capello, 1867 — Asia, Africa
Paraplectanoides Keyserling, 1886 — Australia
Pararaneus Caporiacco, 1940 — Madagascar
Paraverrucosa Mello-Leitão, 1939 — South America
Parawixia F. O. Pickard-Cambridge, 1904 — Mexico, South America, Asia, Papua New Guinea, Central America, Trinidad
Parmatergus Emerit, 1994 — Madagascar
Pasilobus Simon, 1895 — Africa, Asia
Perilla Thorell, 1895 — Myanmar, Vietnam, Malaysia
Pherenice Thorell, 1899 — Cameroon
Phonognatha Simon, 1894 — Australia
Pitharatus Simon, 1895 — Malaysia, Indonesia
Plebs Joseph & Framenau, 2012 — Oceania, Asia
Poecilarcys Simon, 1895 — Tunisia
Poecilopachys Simon, 1895 — Oceania
Poltys C. L. Koch, 1843 — Asia, Africa, Oceania
Popperaneus Cabra-García & Hormiga, 2020 — Brazil, Paraguay
Porcataraneus Mi & Peng, 2011 — India, China
Pozonia Schenkel, 1953 — Caribbean, Paraguay, Mexico, Panama
Prasonica Simon, 1895 — Africa, Asia, Oceania
Prasonicella Grasshoff, 1971 — Madagascar, Seychelles
Pronoides Schenkel, 1936 — Asia
Pronous Keyserling, 1881 — Malaysia, Mexico, Central America, South America, Madagascar
Pseudartonis Simon, 1903 — Africa
Pseudopsyllo Strand, 1916 — Cameroon
Psyllo Thorell, 1899 — Cameroon, Congo
Pycnacantha Blackwall, 1865 — Africa
Rubrepeira Levi, 1992 — Mexico, Brazil
Salsa Framenau & Castanheira, 2022 — Australia, New Caledonia, Papua New Guinea
Scoloderus Simon, 1887 — Belize, North America, Argentina, Caribbean
Sedasta Simon, 1894 — West Africa
Singa C. L. Koch, 1836 — Africa, Asia, North America, Europe
Singafrotypa Benoit, 1962 — Africa
Siwa Grasshoff, 1970 — Asia
Socca Framenau, Castanheira & Vink, 2022 — Australia
Spilasma Simon, 1897 — South America, Honduras
Spinepeira Levi, 1995 — Peru
Spintharidius Simon, 1893 — South America, Cuba
Taczanowskia Keyserling, 1879 — Mexico, South America
Talthybia Thorell, 1898 — China, Myanmar
Tatepeira Levi, 1995 — South America, Honduras
Telaprocera Harmer & Framenau, 2008 — Australia
Testudinaria Taczanowski, 1879 — South America, Panama
Thelacantha Hasselt, 1882 — Madagascar, Asia, Australia
Thorellina Berg, 1899 — Myanmar, Papua New Guinea
Togacantha Dahl, 1914 — Africa
Trichonephila Dahl, 1911 — Africa, Asia, Oceania, North America, South America
Umbonata Grasshoff, 1971 — Tanzania
Ursa Simon, 1895 — Asia, South America, South Africa
Venomius Rossi, Castanheira, Baptista & Framenau, 2023 — Australia
Verrucosa McCook, 1888 — North America, Panama, South America, Australia
Wagneriana F. O. Pickard-Cambridge, 1904 — South America, Central America, Caribbean, North America
Witica O. Pickard-Cambridge, 1895 — Cuba, Mexico, Peru
Wixia O. Pickard-Cambridge, 1882 — Brazil, Guyana, Bolivia
Xylethrus Simon, 1895 — South America, Mexico, Jamaica, Panama
Yaginumia Archer, 1960 — Asia
Zealaranea Court & Forster, 1988 — New Zealand
Zilla C. L. Koch, 1834 — Azerbaijan, India, China
Zygiella F. O. Pickard-Cambridge, 1902 — North America, Asia, Ukraine, South America
| Biology and health sciences | Spiders | Animals |
428502 | https://en.wikipedia.org/wiki/Hand%20washing | Hand washing | Hand washing (or handwashing), also known as hand hygiene, is the act of cleaning one's hands with soap or handwash and water to remove viruses/bacteria/microorganisms, dirt, grease, and other harmful or unwanted substances stuck to the hands. Drying of the washed hands is part of the process as wet and moist hands are more easily recontaminated. If soap and water are unavailable, hand sanitizer that is at least 60% (v/v) alcohol in water can be used as long as hands are not visibly excessively dirty or greasy. Hand hygiene is central to preventing the spread of infectious diseases in home and everyday life settings.
The World Health Organization (WHO) recommends washing hands for at least 20 seconds before and after certain activities. These include the five critical times during the day where washing hands with soap is important to reduce fecal-oral transmission of disease: after using the toilet (for urination, defecation, menstrual hygiene), after cleaning a child's bottom (changing diapers), before feeding a child, before eating and before/after preparing food or handling raw meat, fish, or poultry.
When neither hand washing nor using hand sanitizer is possible, hands can be cleaned with uncontaminated ash and clean water, although the benefits and harms are uncertain for reducing the spread of viral or bacterial infections. However, frequent hand washing can lead to skin damage due to drying of the skin. Moisturizing lotion is often recommended to keep the hands from drying out; dry skin can lead to skin damage which can increase the risk for the transmission of infection.
Steps and duration
The United States Centers for Disease Control and Prevention (CDC) recommends the following steps when washing one's hands for the prevention of transmission of disease:
Wet hands with warm or cold running water. Running water is recommended because standing basins may be contaminated, while the temperature of the water does not seem to make a difference, however some experts suggest warm, tepid water may be superior.
Lather hands by rubbing them with a generous amount of soap, including the backs of hands, between fingers, and under nails. Soap lifts pathogens from the skin, and studies show that people tend to wash their hands more thoroughly when soap is used rather than water alone.
Scrub for at least 20 seconds. Scrubbing creates friction, which helps remove pathogens from skin, and scrubbing for longer periods removes more pathogens.
Rinse well under running water. Rinsing in a basin can recontaminate hands.
Dry with a clean towel or allow to air dry. Wet and moist hands are more easily recontaminated.
The most commonly missed areas are the thumb, the wrist, the areas between the fingers, and under fingernails. Artificial nails and chipped nail polish may harbor microorganisms.
When it is recommended
There are five critical times during the day where washing hands with soap is important to reduce fecal-oral transmission of disease: after using the toilet (for urination, defecation, menstrual hygiene), after cleaning a child's bottom (changing diapers), before feeding a child, before eating and before/after preparing food or handling raw meat, fish, or poultry. Other occasions when correct handwashing technique should be practiced in order to prevent the transmission of disease include before and after treating a cut or wound; after sneezing, coughing, or blowing your nose; after touching animal waste or handling animals; and after touching garbage.
Public health
Health benefits
Hand washing has many significant health benefits, including minimizing the spread of influenza, COVID-19, and other infectious diseases; preventing infectious causes of diarrhea; decreasing respiratory infections;
and reducing infant mortality rate at home birth deliveries. A 2013 study showed that improved hand washing practices may lead to small improvements in the length growth in children under five years of age. In developing countries, childhood mortality rates related to respiratory and diarrheal diseases can be reduced by introducing simple behavioral changes, such as hand washing with soap. This simple action can reduce the rate of mortality from these diseases by almost 50%. Interventions that promote hand washing can reduce diarrhoea episodes by about a third, and this is comparable to providing clean water in low income areas. 48% of reductions in diarrhoea episodes can be associated with hand washing with soap.
Handwashing with soap is the single most effective and inexpensive way to prevent diarrhea and acute respiratory infections (ARI), as automatic behavior performed in homes, schools, and communities worldwide. Pneumonia, a major ARI, is the number one cause of mortality among children under five years old, taking the lives of an estimated 1.8 million children per year. Diarrhea and pneumonia together account for almost 3.5 million child deaths annually. According to UNICEF, turning handwashing with soap before eating and after using the toilet into an ingrained habit can save more lives than any single vaccine or medical intervention, cutting deaths from diarrhea by almost half and deaths from acute respiratory infections by one-quarter. Hand washing is usually integrated with other sanitation interventions as part of water, sanitation, and hygiene (WASH) programmes. Hand washing also protects against impetigo which is transmitted through direct physical contact.
Adverse effects
A small detrimental effect of handwashing is that frequent hand washing can lead to skin damage due to the drying of the skin. A 2012 Danish study found that excessive hand washing can lead to an itchy, flaky skin condition known as contact dermatitis, which is especially common among health-care workers.
Behavior change
In many countries, there is a low rate of hand washing with soap. A study of hand washing in 54 countries in 2015 found that on average, 38.7% of households practiced hand washing with soap.
A 2014 study showed that Saudi Arabia had the highest rate of 97%; the United States near the middle with 77%; and China with the lowest rate of 23%.
Several behavior change methodologies now exist to increase uptake of the behavior of hand washing with soap at the critical times.
Group hand washing for school children at set times of the day is one option in developing countries to engrain hand washing in children's behaviors. The "Essential Health Care Program" implemented by the Department of Education in the Philippines is an example of at scale action to promote children's health and education. Deworming twice a year, supplemented with washing hands daily with soap, brushing teeth daily with fluoride, is at the core of this national program. It has also been successfully implemented in Indonesia.
Substances used
Soap and detergents
Removal of microorganisms from skin is enhanced by the addition of soaps or detergents to water. Soap and detergents are surfactants that kill microorganisms by disorganizing their membrane lipid bilayer and denaturing their proteins. It also emulsifies oils, enabling them to be carried away by running water.
Solid soap
Solid soap, because of its reusable nature, may hold bacteria acquired from previous uses. A small number of studies which have looked at the bacterial transfer from contaminated solid soap have concluded transfer is unlikely as the bacteria are rinsed off with the foam. The CDC still states "liquid soap with hands-free controls for dispensing is preferable".
Antibacterial soap
Antibacterial soaps have been heavily promoted to a health-conscious public. To date, there is no evidence that using recommended antiseptics or disinfectants selects for antibiotic-resistant organisms in nature. However, antibacterial soaps contain common antibacterial agents such as triclosan, which has an extensive list of resistant strains of organisms. So, even if antibiotic resistant strains are not selected for by antibacterial soaps, they might not be as effective as they are marketed to be. Besides the surfactant and skin-protecting agent, the sophisticated formulations may contain acids (acetic acid, ascorbic acid, lactic acid) as pH regulator, antimicrobially active benzoic acid and further skin conditioners (aloe vera, vitamins, menthol, plant extracts).
A 2007 meta-analysis from the University of Oregon School of Public Health indicated that plain soaps are as effective as consumer-grade anti-bacterial soaps containing triclosan in preventing illness and removing bacteria from the hands. Dissenting, a 2011 meta-analysis in the Journal of Food Protection argued that when properly formulated, triclosan can grant a small but detectable improvement, as can chlorhexidine gluconate, iodophor, or povidone.
Warm water
Hot water that is still comfortable for washing hands is not hot enough to kill bacteria. Bacteria grow much faster at body temperature (37 °C). WHO considers warm soapy water to be more effective than cold, soapy water at removing natural oils which hold soils and bacteria. But CDC mentions that warm water causes skin irritations more often and its ecological footprint is more significant. Water temperatures from 4 to 40 °C do not differ significantly regarding removal of microbes. The most important factor is proper scrubbing.
Contrary to popular belief, scientific studies have shown that using warm water has no effect on reducing the microbial load on hands. Using hot water for handwashing can even be regarded as a waste of energy.
Antiseptics (hand sanitizer)
In situations where hand washing with soap is not an option (e.g., when in a public place with no access to wash facilities), a waterless hand sanitizer such as an alcohol hand gel can be used. They can be used in addition to hand washing to minimize risks when caring for "at-risk" groups. To be effective, alcohol hand gels should contain not less than 60%v/v alcohol. Enough hand antiseptic or alcohol rub must be used to thoroughly wet or cover both hands. The front and back of both hands and between and the ends of all fingers must be rubbed for approximately 30 seconds until the liquid, foam or gel is dry. Finger tips must be washed well too, rubbing them in both palms.
A hand sanitizer or hand antiseptic is a non-water-based hand hygiene agent. In the late 1990s and early part of the 21st century, alcohol rub non-water-based hand hygiene agents (also known as alcohol-based hand rubs, antiseptic hand rubs, or hand sanitizers) began to gain popularity. Most are based on isopropyl alcohol or ethanol formulated together with a thickening agent such as Carbomer (polymer of acrylic acid) into a gel, or a humectant such as glycerin into a liquid, or foam for ease of use and to decrease the drying effect of the alcohol. Adding diluted hydrogen peroxide increases further the antimicrobial activity.
Hand sanitizers are most effective against bacteria and less effective against some viruses. Alcohol-based hand sanitizers are almost entirely ineffective against norovirus (or Norwalk) type viruses, the most common cause of contagious gastroenteritis.
US Centers for Disease Control and Prevention recommend hand washing with soap over hand sanitizer rubs, particularly when hands are visibly dirty. The increasing use of these agents is based on their ease of use and rapid killing activity against micro-organisms; however, they should not serve as a replacement for proper hand washing unless soap and water are unavailable. Despite their effectiveness, non-water agents do not cleanse the hands of organic material, but simply disinfect them. It is for this reason that hand sanitizers are not as effective as soap and water at preventing the spread of many pathogens, since the pathogens remain on the hands.
Wipes
Hand washing using hand sanitizing wipes is an alternative during traveling in the absence of soap and water. Alcohol-based hand sanitizer should contain at least 60% alcohol.
Ash or mud
Many people in low-income communities cannot afford soap and use ash or soil instead. The World Health Organization recommended ash or sand as an alternative to soap when soap is not available. Use of ash is common in rural areas of developing countries and has in experiments been shown at least as effective as soap for removing pathogens. However, evidence to support the use of ash to wash hands is of poor quality. It is not clear if washing hands with ash is effective at reducing viral or bacterial spreading compared to washing with mud, not washing, or with washing with water alone. One concern is that if the soil or ash is contaminated with microorganisms it may increase the spread of disease rather than decrease it, however, there is also no clear evidence to determine the level of risk. Like soap, ash is also a disinfecting agent because in contact with water, it forms an alkaline solution.
Technologies and design aspects
Low-cost options when water is scarce
Various low-cost options can be made to facilitate hand washing where tap-water and/or soap is not available e.g. pouring water from a hanging jerrycan or gourd with suitable holes and/or using ash if needed in developing countries.
In situations with limited water supply (such as schools or rural areas in developing countries), there are water-conserving solutions, such as "tippy-taps" and other low-cost options. A tippy-tap is a simple technology using a jug suspended by a rope, and a foot-operated lever to pour a small amount of water over the hands and a bar of soap.
Low-cost hand washing technologies for households may differ from facilities for multiple users. For households, options include tippy taps, bucket/container with tap (such as a Veronica Bucket), conventional tap with/without basin, valve/tap fitted to bottles, bucket and cup, camp sink. Options for multiple users include: adapting household technologies for multiple users, water container fitted to a pipe with multiple taps, water container fitted to a pipe with holes.
Advanced technologies
Several companies around the globe have developed technologies that aim to improve the hand washing process. Among the different inventions, there are eco-friendly devices that use 90% less
water and 60% less soap compared to hand washing under a faucet. Another device uses
light-based rays to detect contaminants on the hands after they have been washed.
Certain environments are especially sensitive to the transmission of pathogenic microorganisms, like health care and food production. Organizations attempting to prevent infection transmission in these environments have started using programmed washing cycles that provide sufficient time for scrubbing the hands with soap and rinsing them with water. Combined with AI-powered software, these technological advancements turn the hand-washing process into digital data, allowing individuals to receive insights and improve their hand hygiene practices.
Drying with towels or hand driers
Effective drying of the hands is an essential part of the hand hygiene process. Therefore, the proper drying of hands after washing should be an integral part of the hand hygiene process in health care.
The World Health Organization (WHO) and the Centers for Disease Control and Prevention (CDC) are clear and straightforward concerning hand hygiene, and recommend paper towels and hand dryers equally. Both have stressed the importance of frequent and thorough hand washing followed by their complete drying as a means to stop the spread of pathogens, like COVID-19. Specifically, the World Health Organization recommends that everyone "frequently clean [their] hands..." and "dry [them] thoroughly by using paper towels or a warm air dryer." The CDC report that, "Both [clean towels or air hand dryers] are effective ways to dry hands."
A study in 2020 found that hand dryers and paper towels were both found to be equally hygienic hand-drying solutions.
However, there is some debate over the most effective form of drying in public toilets. A growing volume of research suggests paper towels are much more hygienic than the electric hand dryers found in many public toilets. A review in 2012 concluded that "From a hygiene standpoint, paper towels are superior to air dryers; therefore, paper towels should be recommended for use in locations in which hygiene is paramount, such as hospitals and clinics."
Jet-air dryers were found to be capable of blowing micro-organisms from the hands and the unit and potentially contaminating other users and the environment up to away. In the same study in 2008 (sponsored by the paper-towel industry the European Tissue Symposium), use of a warm-air hand dryer spread micro-organisms only up to from the dryer, and paper towels showed no significant spread of micro-organisms. No studies have found a correlation to hand dryers and human health, however, making these findings inconsequential.
Accessibility
Making hand washing facilities accessible (inclusive) to everyone is crucial to maintain hand washing behavior. Considerations for accessibility include age, disability, seasonality (with rains and muddiness), location and more. Important aspects for good accessibility include: Placement of the technology, paths, ramps, steps, type of tap, soap placement.
Medical use
Medical hand-washing became mandatory long after Hungarian physician Ignaz Semmelweis discovered its effectiveness (in 1846) in preventing disease in a hospital environment. There are electronic devices that provide feedback to remind hospital staff to wash their hands when they forget. One study has found decreased infection rates with their use.
Method
Medical hand-washing is for a minimum of 15 seconds, using generous amounts of soap and water or gel to lather and rub each part of the hands. Hands should be rubbed together with digits interlocking. If there is debris under fingernails, a bristle brush may be used to remove it. Since pathogens may remain in the water on the hands, it is important to rinse well and wipe dry with a clean towel. After drying, the paper towel should be used to turn off the water (and open any exit door if necessary). This avoids re-contaminating the hands from those surfaces.
The purpose of hand-washing in the health-care setting is to remove pathogenic microorganisms ("germs") and avoid transmitting them. The New England Journal of Medicine reports that a lack of hand-washing remains at unacceptable levels in most medical environments, with large numbers of doctors and nurses routinely forgetting to wash their hands before touching patients, thus transmitting microorganisms. One study showed that proper hand-washing and other simple procedures can decrease the rate of catheter-related bloodstream infections by 66%.
The World Health Organization has published a sheet demonstrating standard hand-washing and hand-rubbing in health-care sectors. The draft guidance of hand hygiene by the organization can also be found at its website for public comment. A relevant review was conducted by Whitby et al. Commercial devices can measure and validate hand hygiene, if demonstration of regulatory compliance is required.
The World Health Organization has "Five Moments" for washing hands:
before patient care
after environmental contact
after exposure to blood/body fluids
before an aseptic task, and
after patient care.
The addition of antiseptic chemicals to soap ("medicated" or "antimicrobial" soaps) confers killing action to a hand-washing agent. Such killing action may be desired before performing surgery or in settings in which antibiotic-resistant organisms are highly prevalent.
To 'scrub' one's hands for a surgical operation, it is necessary to have a tap that can be turned on and off without touching it with the hands, some chlorhexidine or iodine wash, sterile towels for drying the hands after washing, and a sterile brush for scrubbing and another sterile instrument for cleaning under the fingernails. All jewelry should be removed. This procedure requires washing the hands and forearms up to the elbow, usually 2–6 minutes. Long scrub-times (10 minutes) are not necessary. When rinsing, water on the forearms must be prevented from running back to the hands. After hand-washing is completed, the hands are dried with a sterile cloth and a surgical gown is donned.
Effectiveness in healthcare settings
To reduce the spread of pathogens, it is better to wash the hands or use a hand antiseptic before and after tending to a sick person.
For control of staphylococcal infections in hospitals, it has been found that the greatest benefit from hand-cleansing came from the first 20% of washing, and that very little additional benefit was gained when hand cleansing frequency was increased beyond 35%. Washing with plain soap results in more than triple the rate of bacterial infectious disease transmitted to food as compared to washing with antibacterial soap.
Comparing hand-rubbing with alcohol-based solution with hand washing with antibacterial soap for a median time of 30 seconds each showed that the alcohol hand-rubbing reduced bacterial contamination 26% more than the antibacterial soap. But soap and water is more effective than alcohol-based hand rubs for reducing H1N1 influenza A virus and Clostridioides difficile spores from hands.
Interventions to improve hand hygiene in healthcare settings can involve education for staff on hand washing, increasing the availability of alcohol-based hand rub, and written and verbal reminders to staff. There is a need for more research into which of these interventions are most effective in different healthcare settings.
Developing countries
In developing countries, hand washing with soap is recognized as a cost-effective, essential tool for achieving good health, and even good nutrition. However, a lack of reliable water supply, soap or hand washing facilities in people's homes, at schools and the workplace make it a challenge to achieve universal hand washing behaviors. For example, in most of rural Africa hand washing taps close to every private or public toilet are scarce, even though cheap options exist to build hand washing stations. However, low hand washing rates can also be the result of engrained habits rather than due to a lack of soap or water.
Hand washing at a global level has its own indicator within Sustainable Development Goal 6, Target 6.2 which states "By 2030, achieve access to adequate and equitable sanitation and hygiene for all and end open defecation, paying special attention to the needs of women and girls and those in vulnerable situations. The corresponding Indicator 6.2.1 is formulated as follows: "Proportion of population using (a) safely managed sanitation services and (b) a hand-washing facility with soap and water" (see map to the right with data worldwide from 2017)."
Promotion campaigns
The promotion and advocacy of hand washing with soap can influence policy decisions, raise awareness about the benefits of hand washing and lead to long-term behavior change of the population. For this to work effectively, monitoring and evaluation are necessary. A systematic review of 70 studies found that community-based approaches are effective at increasing hand washing in LMICs, while social marketing campaigns are less effective.
One example for hand washing promotion in schools is the "Three Star Approach" by UNICEF that encourages schools to take simple, inexpensive steps to ensure that students wash their hands with soap, among other hygienic requirements. When minimum standards are achieved, schools can move from one to ultimately three stars. Building hand washing stations can be a part of hand washing promotion campaigns that are carried out to reduce diseases and child mortality.
Global Handwashing Day is another example of an awareness-raising campaign that is trying to achieve behavior change.
As a result of the ongoing COVID-19 pandemic, UNICEF promoted the adoption of a hand washing emoji.
Designing hand washing facilities that encourage use can use the following aspects:
Nudges, cues and reminders
Hand washing facilities should be placed at convenient locations to encourage people to use them regularly and at the right times; they should be attractive and well maintained.
Cost effectiveness
Few studies have considered the overall cost effectiveness of hand washing in developing countries in relationship to DALYs averted. However, one review suggests that promoting hand washing with soap is significantly more cost-effective than other water and sanitation interventions.
History
The importance of hand washing for human healthparticularly for people in vulnerable circumstances like mothers who had just given birth or wounded soldiers in hospitalswas first recognized in the mid 19th century by two pioneers of hand hygiene: the Hungarian physician Ignaz Semmelweis who worked in Vienna, Austria and Florence Nightingale, the English "founder of modern nursing". At that time most people still believed that infections were caused by foul odors called miasmas.
In the 1980s, foodborne outbreaks and healthcare-associated infections led the United States Centers for Disease Control and Prevention to more actively promote hand hygiene as an important way to prevent the spread of infection. The outbreak of swine flu in 2009 and the COVID-19 pandemic in 2020 led to increased awareness in many countries of the importance of washing hands with soap to protect oneself from such infectious diseases. For example, posters with "correct hand washing techniques" were hung up next to hand washing sinks in public toilets and in the toilets of office buildings and airports in Germany. Research indicates that the COVID pandemic shifted social norms regarding hand washing, making it more prevalent worldwide.
Society and culture
Moral aspects
The phrase "washing one's hands of" something, means declaring one's unwillingness to take responsibility for the thing or share complicity in it. It originates from the bible passage in Matthew where Pontius Pilate washed his hands of the decision to crucify Jesus Christ, but has become a phrase with a much wider usage in some English communities.
In Shakespeare's Macbeth, Lady Macbeth begins to compulsively wash her hands in an attempt to cleanse an imagined stain, representing her guilty conscience regarding crimes she had committed and induced her husband to commit.
| Biology and health sciences | Hygiene and grooming: General | Health |
428584 | https://en.wikipedia.org/wiki/Starburst%20galaxy | Starburst galaxy | A starburst galaxy is one undergoing an exceptionally high rate of star formation, as compared to the long-term average rate of star formation in the galaxy, or the star formation rate observed in most other galaxies.
For example, the star formation rate of the Milky Way galaxy is approximately 3 M☉/yr, while starburst galaxies can experience star formation rates of 100 M☉/yr or more. In a starburst galaxy, the rate of star formation is so large that the galaxy consumes all of its gas reservoir, from which the stars are forming, on a timescale much shorter than the age of the galaxy. As such, the starburst nature of a galaxy is a phase, and one that typically occupies a brief period of a galaxy's evolution. The majority of starburst galaxies are in the midst of a merger or close encounter with another galaxy. Starburst galaxies include M82, NGC 4038/NGC 4039 (the Antennae Galaxies), and IC 10.
Definition
Starburst galaxies are defined by these three interrelated factors:
The rate at which the galaxy is currently converting gas into stars (the star-formation rate, or SFR).
The available quantity of gas from which stars can be formed.
A comparison of the timescale on which star formation consumes the available gas with the age or rotation period of the galaxy.
Commonly used definitions include:
Continued star-formation where the current SFR would exhaust the available gas reservoir in much less time than the age of the Universe (the Hubble Time).
Continued star-formation where the current SFR would exhaust the available gas reservoir in much less time than the dynamical timescale of the galaxy (perhaps one rotation period in a disk type galaxy).
The current SFR, normalized by the past-averaged SFR, is much greater than unity. This ratio is referred to as the "birthrate parameter".
Triggering mechanisms
Mergers and tidal interactions between gas-rich galaxies play a large role in driving starbursts. Galaxies in the midst of a starburst frequently show tidal tails, an indication of a close encounter with another galaxy, or are in the midst of a merger. Turbulence, along with variations of time and space, cause the dense gas within a galaxy to compress and rapidly increase star formation. The efficiency at which the galaxy forms also increases its SFR . These changes in the rate of star formation also led to variations with depletion time, and power a starburst with its own galactic mechanisms rather than merging with another galaxy. Interactions between galaxies that do not merge can trigger unstable rotation modes, such as the bar instability, which causes gas to be funneled towards the nucleus and ignites bursts of star formation near the galactic nucleus. It has been shown that there is a strong correlation between the lopsidedness of a galaxy and the youth of its stellar population, with more lopsided galaxies having younger central stellar populations. As lopsidedness can be caused by tidal interactions and mergers between galaxies, this result gives further evidence that mergers and tidal interactions can induce central star formation in a galaxy and drive a starburst.
Types
Classifying types of starburst galaxies is difficult since starburst galaxies do not represent a specific type in and of themselves. Starbursts can occur in disk galaxies, and irregular galaxies often exhibit knots of starburst spread throughout the irregular galaxy. Nevertheless, astronomers typically classify starburst galaxies based on their most distinct observational characteristics. Some of the categorizations include:
Blue compact galaxies (BCGs). These galaxies are often low mass, low metallicity, dust-free objects. Because they are dust-free and contain a large number of hot, young stars, they are often blue in optical and ultraviolet colours. It was initially thought that BCGs were genuinely young galaxies in the process of forming their first generation of stars, thus explaining their low metal content. However, old stellar populations have been found in most BCGs, and it is thought that efficient mixing may explain the apparent lack of dust and metals. Most BCGs show signs of recent mergers and/or close interactions. Well-studied BCGs include IZw18 (the most metal poor galaxy known), ESO338-IG04 and Haro11.
Blue compact dwarf galaxies (BCD galaxies) are small compact galaxies.
Green Pea galaxies (GPs) are small compact galaxies resembling primordial starbursts. They were found by citizen scientists taking part in the Galaxy Zoo project.
Luminous infrared galaxies (LIRGs).
Ultra-luminous Infrared Galaxies (ULIRGs). These galaxies are generally extremely dusty objects. The ultraviolet radiation produced by the obscured star-formation is absorbed by the dust and reradiated in the infrared spectrum at wavelengths of around 100 micrometres. This explains the extreme red colours associated with ULIRGs. It is not known for sure that the UV radiation is produced purely by star-formation, and some astronomers believe ULIRGs to be powered (at least in part) by active galactic nuclei (AGN). X-ray observations of many ULIRGs that penetrate the dust suggest that many starburst galaxies are double-cored systems, lending support to the hypothesis that ULIRGs are powered by star-formation triggered by major mergers. Well-studied ULIRGs include Arp 220.
Hyperluminous Infrared galaxies (HLIRGs), sometimes called submillimeter galaxies.
Wolf–Rayet galaxies (WR galaxies), galaxies where a large portion of the bright stars are Wolf–Rayet stars. The Wolf–Rayet phase is a relatively short-lived phase in the life of massive stars, typically 10% of the total life-time of these stars and as such any galaxy is likely to contain few of these. However, because the stars are both luminous and have distinctive spectral features, it is possible to identify these stars in the spectra of entire galaxies and doing so allows good constraints to be placed on the properties of the starbursts in these galaxies.
Ingredients
Firstly, a starburst galaxy must have a large supply of gas available to form stars. The burst itself may be triggered by a close encounter with another galaxy (such as M81/M82), a collision with another galaxy (such as the Antennae), or by another process that forces material into the centre of the galaxy (such as a stellar bar).
The inside of the starburst is quite an extreme environment. The large amounts of gas mean that massive stars are formed. Young, hot stars ionize the gas (mainly hydrogen) around them, creating H II regions. Groups of hot stars are known as OB associations. These stars burn bright and fast, and are quite likely to explode at the end of their lives as supernovae.
After the supernova explosion, the ejected material expands and becomes a supernova remnant. These remnants interact with the surrounding environment within the starburst (the interstellar medium) and can be the site of naturally occurring masers.
Studying nearby starburst galaxies can help us determine the history of galaxy formation and evolution. Large numbers of the most distant galaxies seen, for example, in the Hubble Deep Field are known to be starbursts, but they are too far away to be studied in any detail. Observing nearby examples and exploring their characteristics can give us an idea of what was happening in the early universe as the light we see from these distant galaxies left them when the universe was much younger (see redshift). However, starburst galaxies seem to be quite rare in our local universe, and are more common further away – indicating that there were more of them billions of years ago. All galaxies were closer together then, and therefore more likely to be influenced by each other's gravity. More frequent encounters produced more starbursts as galactic forms evolved with the expanding universe.
Examples
M82 is the archetypal starburst galaxy. Its high level of star formation is due to a close encounter with the nearby spiral M81. Maps of the regions made with radio telescopes show large streams of neutral hydrogen connecting the two galaxies, also as a result of the encounter. Radio images of the central regions of M82 also show a large number of young supernova remnants, left behind when the more massive stars created in the starburst came to the end of their lives. The Antennae is another starburst system, detailed by a Hubble picture, released in 1997.
List of starburst galaxies
Gallery
| Physical sciences | Stellar astronomy | null |
428874 | https://en.wikipedia.org/wiki/Crassulacean%20acid%20metabolism | Crassulacean acid metabolism | Crassulacean acid metabolism, also known as CAM photosynthesis, is a carbon fixation pathway that evolved in some plants as an adaptation to arid conditions that allows a plant to photosynthesize during the day, but only exchange gases at night. In a plant using full CAM, the stomata in the leaves remain shut during the day to reduce evapotranspiration, but they open at night to collect carbon dioxide () and allow it to diffuse into the mesophyll cells. The is stored as four-carbon malic acid in vacuoles at night, and then in the daytime, the malate is transported to chloroplasts where it is converted back to , which is then used during photosynthesis. The pre-collected is concentrated around the enzyme RuBisCO, increasing photosynthetic efficiency. This mechanism of acid metabolism was first discovered in plants of the family Crassulaceae.
Historical background
Observations relating to CAM were first made by de Saussure in 1804 in his Recherches Chimiques sur la Végétation. Benjamin Heyne in 1812 noted that Bryophyllum leaves in India were acidic in the morning and tasteless by afternoon. These observations were studied further and refined by Aubert, E. in 1892 in his Recherches physiologiques sur les plantes grasses and expounded upon by Richards, H. M. 1915 in Acidity and Gas Interchange in Cacti, Carnegie Institution. The term CAM may have been coined by Ranson and Thomas in 1940, but they were not the first to discover this cycle. It was observed by the botanists Ranson and Thomas, in the succulent family Crassulaceae (which includes jade plants and Sedum). The name "Crassulacean acid metabolism" refers to acid metabolism in Crassulaceae, and not the metabolism of "crassulacean acid"; there is no chemical by that name.
Overview: a two-part cycle
CAM is an adaptation for increased efficiency in the use of water, and so is typically found in plants growing in arid conditions. (CAM is found in over 99% of the known 1700 species of Cactaceae and in nearly all of the cacti producing edible fruits.)
During the night
During the night, a plant employing CAM has its stomata open, allowing to enter and be fixed as organic acids by a PEP reaction similar to the pathway. The resulting organic acids are stored in vacuoles for later use, as the Calvin cycle cannot operate without ATP and NADPH, products of light-dependent reactions that do not take place at night.
During the day
During the day, the stomata close to conserve water, and the -storing organic acids are released from the vacuoles of the mesophyll cells. An enzyme in the stroma of chloroplasts releases the , which enters into the Calvin cycle so that photosynthesis may take place.
Benefits
The most important benefit of CAM to the plant is the ability to leave most leaf stomata closed during the day. Plants employing CAM are most common in arid environments, where water is scarce. Being able to keep stomata closed during the hottest and driest part of the day reduces the loss of water through evapotranspiration, allowing such plants to grow in environments that would otherwise be far too dry. Plants using only carbon fixation, for example, lose 97% of the water they take up through the roots to transpiration - a high cost avoided by plants able to employ CAM.
Comparison with metabolism
The pathway bears resemblance to CAM; both act to concentrate around RuBisCO, thereby increasing its efficiency. CAM concentrates it temporally, providing during the day, and not at night, when respiration is the dominant reaction. plants, in contrast, concentrate spatially, with a RuBisCO reaction centre in a "bundle sheath cell" being inundated with . Due to the inactivity required by the CAM mechanism, carbon fixation has a greater efficiency in terms of PGA synthesis.
There are some /CAM intermediate species, such as Peperomia camptotricha, Portulaca oleracea, and Portulaca grandiflora. It was previously thought that the two pathways of photosynthesis in such plants could occur in the same leaves but not in the same cells, and that the two pathways could not couple but only occur side by side. It is now known, however, that in at least some species such as Portulaca oleracea, C4 and CAM photosynthesis are fully integrated within the same cells, and that CAM-generated metabolites are incorporated directly into the C4 cycle.
Biochemistry
Plants with CAM must control storage of and its reduction to branched carbohydrates in space and time.
At low temperatures (frequently at night), plants using CAM open their stomata, molecules diffuse into the spongy mesophyll's intracellular spaces and then into the cytoplasm. Here, they can meet phosphoenolpyruvate (PEP), which is a phosphorylated triose. During this time, the plants are synthesizing a protein called PEP carboxylase kinase (PEP-C kinase), whose expression can be inhibited by high temperatures (frequently at daylight) and the presence of malate. PEP-C kinase phosphorylates its target enzyme PEP carboxylase (PEP-C). Phosphorylation dramatically enhances the enzyme's capability to catalyze the formation of oxaloacetate, which can be subsequently transformed into malate by NAD+ malate dehydrogenase. Malate is then transported via malate shuttles into the vacuole, where it is converted into the storage form malic acid. In contrast to PEP-C kinase, PEP-C is synthesized all the time but almost inhibited at daylight either by dephosphorylation via PEP-C phosphatase or directly by binding malate. The latter is not possible at low temperatures, since malate is efficiently transported into the vacuole, whereas PEP-C kinase readily inverts dephosphorylation.
In daylight, plants using CAM close their guard cells and discharge malate that is subsequently transported into chloroplasts. There, depending on plant species, it is cleaved into pyruvate and either by malic enzyme or by PEP carboxykinase. is then introduced into the Calvin cycle, a coupled and self-recovering enzyme system, which is used to build branched carbohydrates. The by-product pyruvate can be further degraded in the mitochondrial citric acid cycle, thereby providing additional molecules for the Calvin Cycle. Pyruvate can also be used to recover PEP via pyruvate phosphate dikinase, a high-energy step, which requires ATP and an additional phosphate. During the following cool night, PEP is finally exported into the cytoplasm, where it is involved in fixing carbon dioxide via malate.
Use by plants
Plants use CAM to different degrees. Some are "obligate CAM plants", i.e. they use only CAM in photosynthesis, although they vary in the amount of they are able to store as organic acids; they are sometimes divided into "strong CAM" and "weak CAM" plants on this basis. Other plants show "inducible CAM", in which they are able to switch between using either the or mechanism and CAM depending on environmental conditions. Another group of plants employ "CAM-cycling", in which their stomata do not open at night; the plants instead recycle produced by respiration as well as storing some during the day.
Plants showing inducible CAM and CAM-cycling are typically found in conditions where periods of water shortage alternate with periods when water is freely available. Periodic drought – a feature of semi-arid regions – is one cause of water shortage. Plants which grow on trees or rocks (as epiphytes or lithophytes) also experience variations in water availability. Salinity, high light levels and nutrient availability are other factors which have been shown to induce CAM.
Since CAM is an adaptation to arid conditions, plants using CAM often display other xerophytic characters, such as thick, reduced leaves with a low surface-area-to-volume ratio; thick cuticle; and stomata sunken into pits. Some shed their leaves during the dry season; others (the succulents) store water in vacuoles. CAM also causes taste differences: plants may have an increasingly sour taste during the night yet become sweeter-tasting during the day. This is due to malic acid being stored in the vacuoles of the plants' cells during the night and then being used up during the day.
Aquatic CAM
CAM photosynthesis is also found in aquatic species in at least 4 genera, including: Isoetes, Crassula, Littorella, Sagittaria, and possibly Vallisneria, being found in a variety of species e.g. Isoetes howellii, Crassula aquatica.
These plants follow the same nocturnal acid accumulation and daytime deacidification as terrestrial CAM species. However, the reason for CAM in aquatic plants is not due to a lack of available water, but a limited supply of . is limited due to slow diffusion in water, 10000x slower than in air. The problem is especially acute under acid pH, where the only inorganic carbon species present is , with no available bicarbonate or carbonate supply.
Aquatic CAM plants capture carbon at night when it is abundant due to a lack of competition from other photosynthetic organisms. This also results in lowered photorespiration due to less photosynthetically generated oxygen.
Aquatic CAM is most marked in the summer months when there is increased competition for , compared to the winter months. However, in the winter months CAM still has a significant role.
Ecological and taxonomic distribution of CAM-using plants
The majority of plants possessing CAM are either epiphytes (e.g., orchids, bromeliads) or succulent xerophytes (e.g., cacti, cactoid Euphorbias), but CAM is also found in hemiepiphytes (e.g., Clusia); lithophytes (e.g., Sedum, Sempervivum); terrestrial bromeliads; wetland plants (e.g., Isoetes, Crassula (Tillaea), Lobelia); and in one halophyte, Mesembryanthemum crystallinum; one non-succulent terrestrial plant, (Dodonaea viscosa) and one mangrove associate (Sesuvium portulacastrum).
The only trees that can do CAM are in the genus Clusia; species of which are found across Central America, South America and the Caribbean. In Clusia, CAM is found in species that inhabit hotter, drier ecological niches, whereas species living in cooler montane forests tend to be . In addition, some species of Clusia can temporarily switch their photosynthetic physiology from to CAM, a process known as facultative CAM. This allows these trees to benefit from the elevated growth rates of photosynthesis, when water is plentiful, and the drought tolerant nature of CAM, when the dry season occurs.
Plants which are able to switch between different methods of carbon fixation include Portulacaria afra, better known as Dwarf Jade Plant, which normally uses C3 fixation but can use CAM if it is drought-stressed, and Portulaca oleracea, better known as Purslane, which normally uses fixation but is also able to switch to CAM when drought-stressed.
CAM has evolved convergently many times. It occurs in 16,000 species (about 7% of plants), belonging to over 300 genera and around 40 families, but this is thought to be a considerable underestimate. The great majority of plants using CAM are angiosperms (flowering plants) but it is found in ferns, Gnetopsida and in quillworts (relatives of club mosses). Interpretation of the first quillwort genome in 2021 (I. taiwanensis) suggested that its use of CAM was another example of convergent evolution. In Tillandsia CAM evolution has been associated with gene family expansion.
The following list summarizes the taxonomic distribution of plants with CAM:
| Biology and health sciences | Metabolic processes | Biology |
428932 | https://en.wikipedia.org/wiki/Edge%20of%20chaos | Edge of chaos | The edge of chaos is a transition space between order and disorder that is hypothesized to exist within a wide variety of systems. This transition zone is a region of bounded instability that engenders a constant dynamic interplay between order and disorder.
Even though the idea of the edge of chaos is an abstract one, it has many applications in such fields as ecology, business management, psychology, political science, and other domains of the social sciences. Physicists have shown that adaptation to the edge of chaos occurs in almost all systems with feedback.
History
The phrase edge of chaos was coined in the late 1980s by chaos theory physicist Norman Packard. In the next decade, Packard and mathematician Doyne Farmer co-authored many papers on understanding how self-organization and order emerges at the edge of chaos. One of the original catalysts that led to the idea of the edge of chaos were the experiments with cellular automata done by computer scientist Christopher Langton where a transition phenomenon was discovered. The phrase refers to an area in the range of a variable, λ (lambda), which was varied while examining the behaviour of a cellular automaton (CA). As λ varied, the behaviour of the CA went through a phase transition of behaviours. Langton found a small area conducive to produce CAs capable of universal computation. At around the same time physicist James P. Crutchfield and others used the phrase onset of chaos to describe more or less the same concept.
In the sciences in general, the phrase has come to refer to a metaphor that some physical, biological, economic and social systems operate in a region between order and either complete randomness or chaos, where the complexity is maximal.
The generality and significance of the idea, however, has since been called into question by Melanie Mitchell and others. The phrase has also been borrowed by the business community and is sometimes used inappropriately and in contexts that are far from the original scope of the meaning of the term.
Stuart Kauffman has studied mathematical models of evolving systems in which the rate of evolution is maximized near the edge of chaos.
Adaptation
Adaptation plays a vital role for all living organisms and systems. All of them are constantly changing their inner properties to better fit in the current environment. The most important instruments for the adaptation are the self-adjusting parameters inherent for many natural systems. The prominent feature of systems with self-adjusting parameters is an ability to avoid chaos. The name for this phenomenon is "Adaptation to the edge of chaos".
Adaptation to the edge of chaos refers to the idea that many complex adaptive systems (CASs) seem to intuitively evolve toward a regime near the boundary between chaos and order. Physics has shown that edge of chaos is the optimal settings for control of a system. It is also an optional setting that can influence the ability of a physical system to perform primitive functions for computation. In CAS, coevolution generally occurs near the edge of chaos, and a balance should be maintained between flexibility and stability to avoid structural failure. As a response to coping with turbulent environments, CAS bring out flexibility, creativity, agility, anti-fragility, and innovation near the edge of chaos, provided these systems are sufficiently decentralized and non-hierarchical.
Because of the importance of adaptation in many natural systems, adaptation to the edge of the chaos takes a prominent position in many scientific researches. Physicists demonstrated that adaptation to state at the boundary of chaos and order occurs in population of cellular automata rules which optimize the performance evolving with a genetic algorithm. Another example of this phenomenon is the self-organized criticality in avalanche and earthquake models.
The simplest model for chaotic dynamics is the logistic map. Self-adjusting logistic map dynamics exhibit adaptation to the edge of chaos. Theoretical analysis allowed prediction of the location of the narrow parameter regime near the boundary to which the system evolves.
| Mathematics | Dynamical systems | null |
429149 | https://en.wikipedia.org/wiki/Nutria | Nutria | The nutria () or coypu () (Myocastor coypus) is a herbivorous, semiaquatic rodent from South America.
Classified for a long time as the only member of the family Myocastoridae, Myocastor has since been included within Echimyidae, the family of the spiny rats.
The nutria lives in burrows alongside stretches of water and feeds on river plant stems. Originally native to subtropical and temperate South America, it was introduced to North America, Europe and Asia, primarily by fur farmers. Although it is still hunted and trapped for its fur in some regions, its destructive burrowing and feeding habits often bring it into conflict with humans, and it is considered an invasive species in the United States. Nutria also transmit various diseases to humans and animals, mainly through water contamination.
Etymology
The genus name Myocastor derives from the two Ancient Greek words () 'rat, mouse', and () 'beaver'. Therefore, the name Myocastor literally means 'mouse beaver'.
Two names are commonly used in English for Myocastor coypus. The name nutria (from the Spanish word nutria 'otter') is generally used in North America, Asia, and throughout countries of the former Soviet Union; however, in most Spanish-speaking countries, the word nutria refers primarily to the otter. To avoid this ambiguity, the name coypu or coipo (derived from Mapudungun) is used in South America, Britain and other parts of Europe. In France, the nutria is known as a ragondin. In Dutch, it is known as beverrat 'beaver rat'. In German, it is known as Nutria, Biberratte 'beaver rat', or Sumpfbiber 'swamp beaver'. In Italy, instead, the popular name is, as in North America and Asia, nutria, but it is also called castorino 'little beaver', by which its fur is known in Italy. In Swedish, the animal is known as sumpbäver 'marsh/swamp beaver'. In Brazil, the animal is known as ratão-do-banhado 'big swamp rat', nútria, or caxingui (the last from Tupi).
Taxonomy
The nutria was first described by Juan Ignacio Molina in 1782 as Mus coypus, a member of the mouse genus. The genus Myocastor was assigned in 1792 by Robert Kerr. Geoffroy Saint-Hilaire, independently of Kerr, named the species Myopotamus coypus, and it is occasionally referred to by this name.
Four subspecies are generally recognized:
M. c. bonariensis: northern Argentina, Bolivia, Paraguay, Uruguay, southern Brazil (RS, SC, PR, and SP)
M. c. coypus: central Chile, Bolivia
M. c. melanops: Chiloé Island
M. c. santacruzae: Patagonia
M. c. bonariensis, the subspecies present in the northernmost (subtropical) part of the nutria's range, is believed to be the type of nutria most commonly introduced to other continents.
Phylogeny
Comparison of DNA and protein sequences showed that the genus Myocastor is the sister group to the genus Callistomys (painted tree-rats). In turn, these two taxa share evolutionary affinities with other Myocastorini genera: Proechimys and Hoplomys (armored rats) on the one hand, and Thrichomys on the other hand.
Appearance
The nutria somewhat resembles a very large rat, or a beaver with a small, long and skinny hairless tail. Adults are typically in weight, and in body length, with a tail. It is possible for nutria to weigh up to , although adults usually average . Nutria have three sets of fur. The guard hairs on the outer coat are three inches long. They have coarse, darkish brown midlayer fur with soft dense grey under fur, also called the nutria. Three distinguishing features are a white patch on the muzzle, webbed hind feet, and large, bright orange-yellow incisors. They have approximately 20 teeth with four large incisors that grow during the entirety of their lives. The orange discoloration is due to pigment staining from the mineral iron in the tooth enamel. Nutria have prominent four inch long whiskers on each side of their muzzle or cheek area. The mammary glands and teats of female nutria are high on her flanks, to allow their young to feed while the female is in the water. There is no visible distinction between male and female nutria. Both are similar in coloring and weight.
A nutria is often mistaken for a muskrat (Ondatra zibethicus), another widely dispersed, semiaquatic rodent that occupies the same wetland habitats. The muskrat, however, is smaller and more tolerant of cold climates, and has a laterally flattened tail it uses to assist in swimming, whereas the tail of a nutria is round. It can also be mistaken for a small beaver, as beavers and nutria have very similar anatomies and habitats. However, beavers' tails are flat and paddle-like, as opposed to the round tails of nutria.
Life history
Nutria can live up to six years in captivity, but individuals rarely live past three years old in the wild. According to one study, 80% of nutrias die within the first year, and less than 15% of a wild population is over 3 years old. A nutria is considered to have reached old age at 4 years old. Male nutria reach sexual maturity as early as four months, and females as early as three months; however, both can have a prolonged adolescence, up to the age of nine months. Once a female is pregnant, gestation lasts 130 days, and she may give birth to as few as one or as many as 13 offspring. The average nutria reproduction is four offspring. Female nutria will mate within two days after offspring are born. The years of reproduction cycle by litter size. Year one might be large, year two litter size will be smaller and year three the litter size will be another larger size. Females can only produce six litters in her life, rarely seven litters. A female on average will have two litters a year.
Nutria generally line nursery nests with grasses and soft reeds. Baby nutria are precocial, born fully furred and with open eyes; they can eat vegetation and swim with their parents within hours of birth. A female nutria can become pregnant again the day after she gives birth to her young. If timed properly, a female can become pregnant three times within a year. Newborn nutria nurse for seven to eight weeks, after which they leave their mothers. Nutria have been known to be territorial and aggressive when caught or cornered. They will bite and attack humans and dogs when threatened. Nutria are mainly crepuscular or nocturnal, with most activity occurring around dusk and sunset with highest activity around midnight. When food is scarce, nutria will forage during the day. When food is plentiful, nutria will rest and groom during the day.
Distribution
Native to subtropical and temperate South America, its range includes Chile, Argentina, Uruguay, Paraguay and the southern parts of Brazil and Bolivia. It has been introduced to North America, Europe and Asia, primarily by fur ranchers.
The distribution of nutrias outside South America tends to contract or expand with successive cold or mild winters. During cold winters, nutria often suffer frostbite on their tails, leading to infection or death. As a result, populations of nutria often contract and even become locally or regionally extinct as in the Scandinavian countries and such US states as Idaho, Montana, and Nebraska during the 1980s. During mild winters, their ranges tend to expand northward. For example, in recent years, range expansions have been noted in Washington and Oregon, as well as Delaware.
According to the U.S. Geological Survey, nutria were first introduced to the United States in California, in 1899 by William Franklin Frakes. , they had spread to the San Francisco Bay Area, where their digging threatened storm levees, and the California Department of Fish and Wildlife had an active eradication program.
They were first brought to Louisiana in the early 1930s for the fur industry, and the population was kept in check, or at a small population size, because of trapping pressure from the fur traders. The earliest account of nutria spreading freely into Louisiana wetlands from their enclosures was in the early 1940s; a hurricane hit the Louisiana coast for which many people were unprepared, and the storm destroyed the enclosures, enabling the nutria to escape into the wild. According to the Louisiana Department of Wildlife and Fisheries, nutria were also transplanted from Port Arthur, Texas, to the Mississippi River in 1941 and then spread due to a hurricane later that year.
Habitat and feeding
Besides breeding quickly, each nutria consumes large amounts of aquatic vegetation. An individual consumes about 25% of its body weight daily, and feeds year-round. Being one of the world's larger extant rodents, a mature, healthy nutria averages in weight, but they can reach as much as . They eat the base of the above-ground stems of plants, and often dig through soil for roots and rhizomes to eat. Nutria eat parts and whole plants, and go after roots, rhizomes, tubers and black willow tree bark in the winter. Their creation of "eat-outs", areas where a majority of the above- and below-ground biomass has been removed, produces patches in the environment, which in turn disrupts the habitat for other animals and humans dependent on wetlands and marshes. Nutria eat the following plant varieties: cattail, rushes, reeds, arrowheads, flatsedges, and cordgrasses. Commercial crops that nutria also eat are lawn grasses, alfalfa, corn, rice, and sugarcane.
Nutria are found most commonly in freshwater marshes and wetlands, but also inhabit brackish marshes and rarely salt marshes. They either construct their own burrows, or occupy burrows abandoned by beaver, muskrats, or other animals. They are also capable of constructing floating rafts out of vegetation. Nutria live in partially underwater dens. The main chamber is not submerged underground. Nutria are considered to be a species that lives in colonies. One male will share a den with three or four females and their offspring. Nutria use "feeding platforms" which are constructed in the water from cut pieces of vegetation supported by a structure like a log or branches. Muskrat dens and beaver lodges are also often used as feeding platforms.
Commercial use and issues
Farming and the fur trade
Local extinction in their native range due to overharvesting led to the development of nutria fur farms in the late 19th and early 20th centuries. The first farms were in Argentina and then later in Europe, North America, and Asia. These farms have generally not been successful long-term investments, and farmed nutria often are released or escape as operations become unprofitable. The first attempt at nutria farming was in France in the early 1880s, but it was not much of a success. The first efficient and extensive nutria farms were located in South America in the 1920s. The South American farms were very successful, and led to the growth of similar farms in North America and Europe. Nutrias from these farms often escaped, or were deliberately released into the wild to provide a game animal or to remove aquatic vegetation.
Nutria were introduced to the Louisiana ecosystem in the 1930s, when they escaped from fur farms that had imported them from South America. Nutria were released into the wild by at least one Louisiana nutria farmer in 1933 and these releases were followed by E. A. McIlhenny who released his entire stock in 1945 on Avery Island. In 1940, some of the nutria escaped during a hurricane and quickly populated coastal marshes, inland swamps, and other wetland areas. From Louisiana, nutria have spread across the Southern United States, wreaking havoc on marshlands.
Following a decline in demand for nutria fur, nutria have since become pests in many areas, destroying aquatic vegetation, marshes, and irrigation systems, and chewing through man-made items such as tires and wooden house panelling in Louisiana, eroding river banks, and displacing native animals. Damage in Louisiana has been sufficiently severe since the 1950s to warrant legislative attention; in 1958, the first bounty was placed on nutria, though this effort was not funded. By the early 2000s, the Coastwide Nutria Control Program was established, which began paying bounties for nutria killed in 2002. In the Chesapeake Bay region in Maryland, where they were introduced in the 1940s, nutria are believed to have destroyed of marshland in the Blackwater National Wildlife Refuge. In response, by 2003, a multimillion-dollar eradication program was underway.
In the United Kingdom, nutria were introduced to East Anglia, for fur, in 1929; many escaped and damaged the drainage works, and a concerted programme by MAFF eradicated them by 1989.
Food products
A small number of game meat websites on the internet sell nutria meat for consumption. As of 2016, at least one Moscow restaurant serves nutria meat dishes. In 1997 and 1998, Louisiana attempted to encourage the public to consume nutria meat. Nutria meat is leaner with a lower fat content and lower in cholesterol compared to ground beef. In an effort to encourage Louisianians to eat nutria, several recipes were distributed to locals and published on the internet. People in poor and rural Louisiana have trapped and consumed nutria meat for decades.
Marsh Dog, a US company based in Baton Rouge, Louisiana, received a grant from the Barataria-Terrebonne National Estuary Program to establish a company that uses nutria meat for dog food products. In 2012, the Louisiana Wildlife Federation recognized Marsh Dog with "Business Conservationist of the Year" award for finding a use for this eco-sustainable protein. A claimed environmentally sound solution is the use of nutria meat to make dog food treats.
In Kyrgyzstan and Uzbekistan, nutria (Russian and local languages Нутрия) are farmed on private plots and sold in local markets as a poor man's meat. As of 2016, however, the meat is used successfully in Moscow restaurant Krasnodar Bistro, as part of the growing Russian localvore movement and as a 'foodie' craze. It appears on the menu as a burger, hotdog, dumplings, or wrapped in cabbage leaves, with the flavour being somewhere between turkey and pork.
Ecological impacts
Herbivory and habitat degradation
Nutria herbivory "severely reduces overall wetland biomass and can lead to the conversion of wetland to open water." Unlike other common disturbances in marshlands, such as fire and tropical storms, which are a once- or few-times-a-year occurrence, nutria feed year round, so their effects on the marsh are constant. Also, nutria are typically more destructive in the winter than in the growing season, due largely to the scarcity of above-ground vegetation; as nutria search for food, they dig up root networks and rhizomes for food.
While nutria are the most common herbivores in Louisiana marshes, they are not the only ones. Feral hogs, also known as wild boars (Sus scrofa), swamp rabbits (Sylvilagus aquaticus), and muskrats (Ondatra zibethicus) are less common, but feral hogs are increasing in number in Louisiana wetlands. On plots open to nutria herbivory, 40% less vegetation was found than in plots guarded against nutria by fences. This number may seem insignificant, and indeed herbivory alone is not a serious cause of land loss, but when herbivory was combined with an additional disturbance, such as fire, single vegetation removal, or double vegetation removal to simulate a tropical storm, the effect of the disturbances on the vegetation were greatly amplified.
As different factors were added together, they resulted in less overall vegetation. Adding fertilizer to open plots did not promote plant growth; instead, nutria fed more in the fertilized areas. Increasing fertilizer inputs in marshes only increases nutria biomass instead of the intended vegetation, therefore increasing nutrient input is not recommended. The problem became so serious that Sheriff Harry Lee of Jefferson Parish used SWAT sharpshooters against the animals.
Wetlands in general are a valuable resource both economically and environmentally. For instance, the U.S. Fish and Wildlife Service determined wetlands covered only 5% of the land surface of the contiguous 48 United States, but they support 31% of the nation's plant species. These very biodiverse systems provide resources, shelter, nesting sites, and resting sites (particularly Louisiana's coastal wetlands such as Grand Isle for migratory birds) to a wide array of wildlife. Human users also receive many benefits from wetlands, such as cleaner water, storm surge protection, oil and gas resources (especially on the Gulf Coast), reduced flooding, and chemical and biological waste reduction, to name a few. In Louisiana, rapid wetland loss occurs due to a variety of reasons; this state loses an estimated area about the size of a football field every hour.
In 1998, the Louisiana Department of Wildlife and Fisheries (LDWF) conducted the first Louisiana coast-wide survey, which was funded by the Coastal Wetlands Planning, Protection, and Restoration Act and titled the Nutria Harvest and Wetland Demonstration Program, to evaluate the condition of the marshlands. The survey revealed through aerial surveys of transects that herbivory damage to wetlands totaled roughly . The next year, LDWF performed the same survey and found the area damaged by herbivory increased to about . The LDWF has determined the wetlands affected by nutria decreased from an estimated minimum of of Louisiana wetlands in 2002–2003 season to about during the 2010–2011 season. The LDWF stresses that coastal wetland restoration projects will be greatly hindered without effective, sustainable nutria population control.
Pathogenic and viral reservoirs of zoonotic diseases
In addition to direct environmental damage, nutria are the host for a roundworm nematode parasite () that can infect the skin of humans, causing dermatitis similar to strongyloidiasis. The condition is also called "nutria itch". Other parasites they can host are tapeworms, liver flukes, and blood flukes. Waterbody contamination by nutria occurs through urine and feces. Nutria also host fleas, ticks and chewing louse. They can carry several zoonotic diseases (diseases transmitted from animals to humans). They are reservoirs for salmonellosis, encephalomyocarditis virus, chlamydia psittaci and antibiotic resistant bacteria, Aeromonas spp. Other zoonotic disease of concern they are host reservoirs for are mycobacterium tuberculosis, septicemia, toxoplasmosis, and rickettsiosis. According to the CDC, nutria carry two out of eight diseases of concern for the United States, rabies and salmonellosis. Nutria are considered a global alien species and have potential to spread disease to livestock and humans. Nutria are found on every continent except Australia and Antarctica. Native to the southern hemisphere and spreading globally requires preventive monitoring for zoonotic disease transmission. Nutria immigration is monitored for habitat destruction of wetlands, farmlands, marshes and is measured in habitat loss in acres. Increased local awareness of viral, bacterial and parasitic transmission from nutria to humans and livestock will be of greater importance as climate change progresses.
Control efforts
As a global alien species, nutria are monitored and managed throughout the world. Many countries have attempted eradication efforts with varying degrees of success.
Nutria are predicted to expand their range northward over the next century as global temperatures increase.
European Union
This species is included since 2016 in the EU list of Invasive Alien Species of Union concern (the Union list). This implies that this species cannot be imported, bred, transported, commercialized, or intentionally released into the environment in the whole of the European Union.
Ireland
A nutria was first sighted in the wild in Ireland in 2010. Some nutria escaped from a pet farm in Cork City in 2015 and began breeding on the outskirts of the city. Ten were trapped on the Curraheen River in 2017, but the rodents continued to spread, reaching Dublin via the Royal Canal in 2019. Animals were found along the River Mulkear in 2015. The National Biodiversity Data Centre issued a species alert in 2017, saying that nutria "[have] the potential to be a high impact invasive species in Ireland. […] This species is listed as among 100 of the worst invasive species in Europe."
Great Britain
In the UK, nutria escaped from fur farms and were reported in the wild as early as 1932. There were three unsuccessful attempts to control nutria in East Anglia between 1943 and 1944. Nutria population and range increased, causing damage to agriculture in the 1950s. During the 1960s, a grant was awarded to rabbit clearance societies that included nutria. This control allowed for the removal of 97,000 nutria in 1961 and 1962. From 1962 to 1965, 12 trappers were hired to eradicate as many nutria as possible confining the remaining population within the Norfolk Broads. The campaign used live traps allowing non-target species to be released while any nutria caught were shot. Combined with cold winters in 1962 to 1963, almost 40,500 nutria were removed from the population. Although nutria populations were greatly reduced after the 1962–1965 campaign ended, the population increased until another eradication campaign began in 1981. This campaign succeeded in fully eradicating nutria in Great Britain. The trapping areas were broken into eight sectors leaving no area uncontrolled. The 24 trappers were offered an incentive for early completion of the 10-year campaign. In 1989 nutria were assumed eradicated, as only three males were found between 1987 and 1989.
Japan
Nutria were introduced to Japan in 1939. They were imported from France during World War II to support food shortages as well as the fur trade. After the war in 1950, many nutria were released en masse or escaped, and became one of Japan's worst invasive species, damaging river banks, rice fields and other valuable crops. In 1963 an eradication program was started to remove nutria but has shown little to no success. Nutria are still present in Japan and there is a restriction on importing, transporting and obtaining nutria per the Invasive Alien Species Act established in 2004.
New Zealand
Nutria are classed as a "prohibited new organism" under New Zealand's Hazardous Substances and New Organisms Act 1996, preventing it from being imported into the country.
United States
Atlantic Coast
An eradication program on the Delmarva Peninsula, between Chesapeake Bay and the Atlantic Coast, where nutria once numbered in the tens of thousands and had destroyed thousands of hectares of marshland, had nearly succeeded by 2012. In September 2022 government officials announced that nutria have been completely eradicated on the Maryland Eastern Shore.
California
The first records of nutria invading California dates from the 1940s and 1950s, when the species was found in the agriculture-rich Central Valley and the south coast of the state, but by the 1970s the animals had been extirpated statewide. They were found again in Merced County in 2017, on the edge of the San Joaquin River Delta. State officials are concerned that they will harm infrastructure that sends water to San Joaquin Valley farms and urban areas. In 2019, the California Department of Fish and Wildlife (CDFW) received nearly $2 million in Governor Gavin Newsom's first budget, and an additional $8.5 million via the Delta Conservancy (a state agency focused on the Delta) to be spent over the course of three years. The state has adopted an eradication campaign based on the successful effort in the Chesapeake Bay, including strategies such as the "Judas nutria" (in which individualized nutria are caught, sterilized, fitted with radio collars, and released, whereupon they can be tracked by hunters as they return to their colonies) and the use of trained dogs. The state has also reversed a prior "no-hunting" policy, although hunting the animals does require a license. California has a restriction on importation and transportation without a permit. If nutria are found or captured in the state of California, local authorities must be notified right away and the nutria cannot be released. Licensed hunters in the state of California may hunt nutria as a non-game animal. Bounty programs are not advised in California due to native species of muskrat and beaver being misidentified.
Louisiana
The Louisiana Coastwide Nutria Control Program provides incentives for harvesting nutria. Starting in 2002, the Louisiana Department of Wildlife and Fisheries (LDWF) performed aerial surveys just as they had done for the Nutria Harvest and Wetland Demonstration Program, except under a different program title. Under the Coastwide Nutria Control Program, which also receives funds from CWPPRA, 308,160 nutria were harvested the first year (2002–2003), revealing damage and totaling $1,232,640 in incentive payouts to those legally participating in the program. Essentially, once a person receives a license to hunt or trap nutria, that person is allowed to capture an unlimited number. When a nutria is captured, the tail is cut off and turned in to a Coastal Environments Inc. (CEI) official at an approved site. As of 2019, each nutria tail is worth $6, which is an increase from $4 before the 2006–2007 season. Nutria harvesting increased drastically during the 2009–2010 season, with 445,963 nutria tails turned in worth $2,229,815 in incentive payments. Each CEI official keeps record of how many tails have been turned in by each individual per parish, the method used in capture of the nutria, and the location of capture. All of this information is transferred to a database to calculate the density of nutria across the Louisiana coast, and the LDWF combines these data with the results from the aerial surveys to determine the number of nutria remaining in the marshes and the amount of damage they are inflicting on the ecosystem.
Another program executed by LDWF involves creating a market of nutria meat for human consumption, though it is still trying to gain public notice. Nutria is a very lean, protein-rich meat, low in fat and cholesterol with the taste, texture, and appearance of rabbit or dark turkey meat. Few pathogens are associated with the meat, but proper heating when cooking should kill them. The quality of the meat and the minimal harmful microorganisms associated with it make nutria meat an "excellent food product for export markets".
Several desirable control methods are ineffective for various reasons. Zinc phosphide is the only rodenticide registered to control nutria, but it is expensive, remains toxic for months, detoxifies in high humidity and rain, and requires construction of expensive floating rafts for placement of the chemical. It is not yet sure how many nontarget species are susceptible to zinc phosphide, but birds and rabbits have been known to die from ingestion. Therefore, this chemical is rarely used, especially not in large-scale projects. Other potential chemical pesticides would be required by the US Environmental Protection Agency to undergo vigorous testing before being acceptable to use on nutria. The LDWF has estimated costs for new chemicals to be $300,000 for laboratory, chemistry, and field studies, and $500,000 for a mandatory Environmental Impact Statement. Contraception is not a common form of control, but is preferred by some wildlife managers. It is also expensive to operate– an estimated $6 million annually to drop bait laced with birth-control chemicals. Testing of other potential contraceptives would take about five to eight years and $10 million, with no guarantee of FDA approval. Also, an intensive environmental assessment would have to be completed to determine whether any non-target organisms were affected by the contraceptive chemicals. Neither of these control methods is likely to be used in the near future.
In Louisiana, a claimed environmentally sound solution is the killing of nutria to make dog food treats.
Gallery
| Biology and health sciences | Rodents | null |
429167 | https://en.wikipedia.org/wiki/Shaku%20%28unit%29 | Shaku (unit) | or Japanese foot is a Japanese unit of length derived (but varying) from the Chinese , originally based upon the distance measured by a human hand from the tip of the thumb to the tip of the forefinger (compare span). Traditionally, the length varied by location or use, but it is now standardized as 10/33 m, or approximately .
Etymology in English
entered English in the early 18th century, a romanization of the Japanese Go-on reading of the character for .
Use in Japan
The had been standardized as since 1891. This means that there are about 3.3 () to one meter.
This definition was established by Meiji government law; until then, even though the unit was given the same name, its length varied depending on the era. At the same time, other units were established based on shaku.
English：1Shaku = 10Cun = 100bu
Japanese：1尺 = 10寸 = 100分
The use of the unit for official purposes in Japan was banned on March 31, 1966, although it is still used in traditional Japanese carpentry and some other fields, such as kimono construction. The traditional Japanese bamboo flute known as the ( and ) derives its name from its length of one and eight . Similarly, the remains in use in the Japanese lumber trade. In the Japanese construction industry, the standard sizes of drywall, plywood, and other sheet goods are based on , with the most common width being three (rounded up to ).
In Japanese media parlance, refers to screen time: the amount of time someone or something is shown on screen (similar to the English "footage").
History
Traditionally, the actual length of the varied over time, location, and use. By the early 19th century, the was largely within the range of , but a longer value of the (also known as the ) was also known, and was 1.17 times longer than the present value ().
Carpenter's unit and tailor's unit
Another variant was used for measuring cloth, which measured meters (), and was known as the , as baleen (whale whiskers) were used as cloth rulers.
To distinguish the two variants of , the general unit was known as the . The Shōsōin treasure house in Nara preserves some antique ivory one- rulers, known as the .
Derived units
Length
Just as with the Chinese unit, the is divided into ten smaller units, known as in Japanese, and ten together form a larger unit known in Japanese as a . The Japanese also had a third derived unit, the , equal to six ; this was used extensively in traditional Japanese architecture as the distance between supporting pillars in Buddhist temples and Shinto shrines.
Volume
Ten cubic comprised a , reckoned as the amount of rice necessary to sustain a peasant for a year.
Outside Japan
The Japanese also forms the basis of the modern Taiwanese foot.
In 1909, the Korean Empire adopted the Japanese definition of the as that of the ().
| Physical sciences | East Asian | Basics and measurement |
429191 | https://en.wikipedia.org/wiki/Gamut | Gamut | In color reproduction and colorimetry, a gamut, or color gamut , is a convex set containing the colors that can be accurately represented, i.e. reproduced by an output device (e.g. printer or display) or measured by an input device (e.g. camera or visual system). Devices with a larger gamut can represent more colors. Similarly, gamut may also refer to the colors within a defined color space, which is not linked to a specific device. A trichromatic gamut is often visualized as a color triangle. A less common usage defines gamut as the subset of colors contained within an image, scene or video.
Introduction
The term gamut was adopted from the field of music, where the medieval Latin expression "gamma ut" meant the lowest tone of the G scale and, in time, came to imply the entire range of musical notes of which musical melodies are composed. Shakespeare's use of the term in The Taming of the Shrew is sometimes attributed to the author / musician Thomas Morley. In the 1850s, the term was applied to a range of colors or hue, for example by Thomas de Quincey, who wrote "Porphyry, I have heard, runs through as large a gamut of hues as marble."
The gamut of a device or process is that portion of the color space that can be represented, or reproduced. Generally, the color gamut is specified in the hue–saturation plane, as a system can usually produce colors over a wide intensity range within its color gamut; for a subtractive color system (such as used in printing), the range of intensity available in the system is for the most part meaningless without considering system-specific properties (such as the illumination of the ink).
Device gamuts must use real primaries (those that can be represented by a physical spectral power distribution) and therefore are always incomplete (smaller than the human visual gamut). No gamut defined by a finite number of primaries can represent the entire human visual gamut. Three primaries are necessary for representing an approximation of the human visual gamut. More primaries can be used to increase the size of the gamut. For example, while painting with red, yellow and blue pigments is sufficient for modeling color vision, adding further pigments (e.g. orange or green) can increase the size of the gamut, allowing the reproduction of more saturated colors.
While processing a digital image, the most convenient color model used is the RGB model. Printing the image requires transforming the image from the original RGB color model to the printer's CMYK color model. During this process, the colors from the RGB model which are out of gamut must be somehow converted to approximate values within the CMYK model. Simply trimming only the colors which are out of gamut to the closest colors in the destination space would burn the image. There are several algorithms approximating this transformation, but none of them can be truly perfect, since those colors are simply out of the target device's capabilities. This is why identifying the colors in an image that are out of gamut in the target color space as soon as possible during processing is critical for the quality of the final product. It is also important to remember that there are colors inside the CMYK gamut that are outside the most commonly used RGB color spaces, such as sRGB and Adobe RGB.
Color management
Color management is the process of ensuring consistent and accurate colors across devices with different gamuts. Color management handles the transformations between color gamuts and canonical color spaces to ensure that colors are represented equally on different devices. A device's gamut is defined by a color profile, usually the ICC profile, which relates the gamut to a standardized color space and allows for calibration of the device. Transforming from one gamut to a smaller gamut loses information as out-of-gamut colors are projected on to the smaller gamut and transforming back to the larger gamut does not regain this lost information.
Colorimetry
Colorimetry is the measurement of color, generally in a way that mimics human color perception. Input devices such as digital cameras or scanners are made to mimic trichromatic human color perception and are based on three sensors elements with different spectral sensitivities, ideally aligned approximately with the spectral sensitivities of human photopsins. In this sense, they have a similar gamut to the human visual system. However, most of these devices violate the Luther condition and are not intended to be truly colorimetric, with the exception of tristimulus colorimeters. Higher-dimension input devices, such as multispectral imagers, hyperspectral imagers or spectrometers, capture color at a much larger gamut, dimensionally, than the human visual gamut. To be perceived by humans, the images must first be down-dimensionalized and treated with false color.
Visual gamut
The extent of color that can be detected by the average human, approximated by the standard observer, is the human visual gamut. The visual gamut is usually visualized in the CIE 1931 chromaticity diagram, where the spectral locus (curved edge) represents the monochromatic (single-wavelength) or spectral colors. As the device you are using to view the diagram has a smaller gamut than the visual gamut, the colors that are out-of-gamut are reproduced as colors inside the display's gamut. Device gamuts are generally depicted in reference to the visual gamut. The standard observer represents a typical human, but colorblindness leads to a reduced visual gamut.
Color reproduction
Visualization of gamuts
The gamut of reflective colors in nature has a similar, though more rounded, shape. An object that reflects only a narrow band of wavelengths will have a color close to the edge of the CIE diagram, but it will have a very low luminosity at the same time. At higher luminosities, the accessible area in the CIE diagram becomes smaller and smaller, up to a single point of white, where all wavelengths are reflected exactly 100 percent; the exact coordinates of white are determined by the color of the light source.
Limitations of color representation
Surfaces
In the beginning of the 20th century, industrial demands for a controllable way to describe colors and the new possibility to measure light spectra initiated intense research on mathematical descriptions of colors.
The idea of optimal colors was introduced by the Baltic German chemist Wilhelm Ostwald. Erwin Schrödinger showed in his 1919 article (Theory of Pigments with Highest Luminosity) that the most-saturated colors that can be created with a given total reflectivity are generated by surfaces having either zero or full reflectance at any given wavelength, and the reflectivity spectrum must have at most two transitions between zero and full.
Thus two types of "optimal color" spectra are possible: Either the transition goes from zero at both ends of the spectrum to one in the middle, as shown in the image at right, or it goes from one at the ends to zero in the middle. The first type produces colors that are similar to the spectral colors and follow roughly the horseshoe-shaped portion of the CIE xy chromaticity diagram, but are generally less saturated. The second type produces colors that are similar to (but generally less saturated than) the colors on the straight line in the CIE xy chromaticity diagram, leading to magenta-like colors.
Schrödinger's work was further developed by David MacAdam and . MacAdam was the first person to calculate precise coordinates of selected points on the boundary of the optimal color solid in the CIE 1931 color space for lightness levels from Y = 10 to 95 in steps of 10 units. This enabled him to draw the optimal color solid at an acceptable degree of precision. Because of his achievement, the boundary of the optimal color solid is called the MacAdam limit (1935).
In 1980, Michael R. Pointer published a maximum gamut for real surfaces with diffuse reflection using 4089 samples, (surfaces with specular reflection ("glossy") can fall outside of this gamut). Originally called a "Munsell Color Cascade", the limits are more commonly called Pointer's Gamut after his work. This gamut remains important as a reference for color reproduction, having been updated by newer methods in ISO 12640-3 Annex B.
On modern computers, it is possible to calculate an optimal color solid with great precision in seconds. The MacAdam limit, on which the most saturated (or "optimal") colors reside, shows that colors that are near monochromatic colors can only be achieved at very low luminance levels, except for yellows, because a mixture of the wavelengths from the long straight-line portion of the spectral locus between green and red will combine to make a color very close to a monochromatic yellow.
Light sources
Light sources used as primaries in an additive color reproduction system need to be bright, so they are generally not close to monochromatic. That is, the color gamut of most variable-color light sources can be understood as a result of difficulties producing pure monochromatic (single wavelength) light. The best technological source of monochromatic light is the laser, which can be rather expensive and impractical for many systems. However, as optoelectronic technology matures, single-longitudinal-mode diode lasers are becoming less expensive, and many applications can already profit from this; such as Raman spectroscopy, holography, biomedical research, fluorescence, reprographics, interferometry, semiconductor inspection, remote detection, optical data storage, image recording, spectral analysis, printing, point-to-point free-space communications, and fiber optic communications.
Systems that use additive color processes usually have a color gamut which is roughly a convex polygon in the hue-saturation plane. The vertices of the polygon are the most saturated colors the system can produce. In subtractive color systems, the color gamut is more often an irregular region.
Comparison of various systems
Following is a list of representative color systems more-or-less ordered from large to small color gamut:
A laser video projector uses three lasers to produce the broadest gamut available in practical display equipment today, derived from the fact that lasers produce truly monochromatic primaries. The systems work either by scanning the entire picture a dot at a time and modulating the laser directly at high frequency, much like the electron beams in a cathode-ray tube (CRT), or by optically spreading and then modulating the laser and scanning a line at a time, the line itself being modulated in much the same way as in a DLP projector. Lasers can also be used as a light source for a DLP projector. More than three lasers can be combined to increase the gamut range, a technique sometimes used in holography.
Digital Light Processing or DLP technology is a trademarked technology from Texas Instruments. The DLP chip contains a rectangular array of up to 2 million hinge-mounted microscopic mirrors. Each of the micromirrors measures less than one-fifth the width of a human hair. A DLP chip's micromirror tilts either toward the light source in a DLP projection system (ON) or away from it (OFF). This creates a light or dark pixel on the projection surface. Current DLP projectors use a quickly rotating wheel with transparent colored "pie slices" to present each color frame successively. One rotation shows the complete image.
Photographic film can reproduce a larger color gamut than typical television, computer, or home video systems.
CRT and similar video displays have a roughly triangular color gamut which covers a significant portion of the visible color space. In CRTs, the limitations are due to the phosphors in the screen which produce red, green, and blue light.
Liquid crystal display (LCD) screens filter the light emitted by a backlight. The gamut of an LCD screen is therefore limited to the emitted spectrum of the backlight. Typical LCD screens use cold-cathode fluorescent bulbs (CCFLs) for backlights. LCD Screens with certain LED or wide-gamut CCFL backlights yield a more comprehensive gamut than CRTs. However, some LCD technologies vary the color presented by viewing angle. In Plane Switching or Patterned vertical alignment screens have a wider span of colors than Twisted Nematic.
Television normally uses a CRT, LCD, LED or plasma display, but does not take full advantage of its color display properties, due to the limitations of broadcasting. The common color profile for TV is based on ITU standard Rec. 601. HDTV is less restrictive and uses a slightly improved color profile based on ITU standard Rec. 709. Still somewhat less than, for example, computer displays using the same display technology. This is due to the use of a limited subset of RGB in broadcasting (values from 16-235), versus full RGB in computer displays, where all bits from 0 to 255 are used.
Paint mixing, both artistic and for commercial applications, achieves a reasonably large color gamut by starting with a larger palette than the red, green, and blue of CRTs or cyan, magenta, and yellow of printing. Paint may reproduce some highly saturated colors that cannot be reproduced well by CRTs (particularly violet), but overall the color gamut is smaller.
Printing typically uses the CMYK color space (cyan, magenta, yellow, and black). Very few printing processes do not include black; however, those processes (with the exception of dye-sublimation printers) are poor at representing low saturation, low intensity colors. Efforts have been made to expand the gamut of the printing process by adding inks of non-primary colors; these are typically orange and green (see Hexachrome) or light cyan and light magenta (see CcMmYK color model). Spot color inks of a very specific color are also sometimes used.
A monochrome display's color gamut is a one-dimensional curve in color space.
Wide color gamut
The Ultra HD Forum defines a wide color gamut (WCG) as a color gamut wider than that of BT.709 (Rec. 709). Color spaces with WCGs include:
Rec. 2020 – ITU-R Recommendation for UHDTV
Rec. 2100 – ITU-R Recommendation for HDR-TV (same chromaticity of color primaries and white point as Rec. 2020)
DCI-P3
Adobe RGB color space
DxO Wide Gamut
Extended-gamut printing
The print gamut achieved by using cyan, magenta, yellow, and black inks is sometimes a limitation, for example when printing colors of corporate logos. Therefore, some methods of color printing use additional ink colors to achieve a larger gamut. For example, some use green, orange, and violet inks to increase the achievable saturation of hues near those. These method are variously called heptatone color printing, extended gamut printing, and 7-color printing, etc.
| Physical sciences | Basics | Physics |
15361791 | https://en.wikipedia.org/wiki/Oscilloscope | Oscilloscope | An oscilloscope, (formerly known as an
oscillograph), (informally scope or O-scope) is a type of electronic test instrument that graphically displays varying voltages of one or more signals as a function of time. Their main purpose is capturing information on electrical signals for debugging, analysis, or characterization. The displayed waveform can then be analyzed for properties such as amplitude, frequency, rise time, time interval, distortion, and others. Originally, calculation of these values required manually measuring the waveform against the scales built into the screen of the instrument. Modern digital instruments may calculate and display these properties directly.
Oscilloscopes are used in the sciences, engineering, biomedical, automotive and the telecommunications industry. General-purpose instruments are used for maintenance of electronic equipment and laboratory work. Special-purpose oscilloscopes may be used to analyze an automotive ignition system or to display the waveform of the heartbeat as an electrocardiogram, for instance.
History
Early high-speed visualisations of electrical voltages were made with an electro-mechanical oscillograph, invented by André Blondel in 1893. These gave valuable insights into high speed voltage changes, but had a frequency response in single kHz, and were superseded by the oscilloscope which used a cathode-ray tube (CRT) as its display element.
The Braun tube, forerunner of the CRT, was known in 1897, and in 1899 Jonathan Zenneck equipped it with beam-forming plates and a magnetic field for deflecting the trace, and this formed the basis of the CRT. Early CRTs had been applied experimentally to laboratory measurements as early as the 1920s, but suffered from poor stability of the vacuum and the cathode emitters. V. K. Zworykin described a permanently sealed, high-vacuum CRT with a thermionic emitter in 1931. This stable and reproducible component allowed General Radio to manufacture an oscilloscope that was usable outside a laboratory setting.
After World War II surplus electronic parts became the basis for the revival of Heathkit Corporation, and a $50 oscilloscope kit made from such parts proved its premiere market success.
Features and uses
An analog oscilloscope is typically divided into four sections: the display, vertical controls, horizontal controls and trigger controls. The display is usually a CRT with horizontal and vertical reference lines called the graticule. CRT displays also have controls for focus, intensity, and beam finder.
The vertical section controls the amplitude of the displayed signal. This section has a volts-per-division (Volts/Div) selector knob, an AC/DC/Ground selector switch, and the vertical (primary) input for the instrument. Additionally, this section is typically equipped with the vertical beam position knob.
The horizontal section controls the time base or sweep of the instrument. The primary control is the Seconds-per-Division (Sec/Div) selector switch. Also included is a horizontal input for plotting dual X-Y axis signals. The horizontal beam position knob is generally located in this section.
The trigger section controls the start event of the sweep. The trigger can be set to automatically restart after each sweep or can be configured to respond to an internal or external event. The principal controls of this section are the source and coupling selector switches, and an external trigger input (EXT Input) and level adjustment.
In addition to the basic instrument, most oscilloscopes are supplied with a probe. The probe connects to any input on the instrument and typically has a resistor of ten times the oscilloscope's input impedance. This results in a 0.1 (‑10×) attenuation factor; this helps to isolate the capacitive load presented by the probe cable from the signal being measured. Some probes have a switch allowing the operator to bypass the resistor when appropriate.
Size and portability
Most modern oscilloscopes are lightweight, portable instruments compact enough for a single person to carry. In addition to portable units, the market offers a number of miniature battery-powered instruments for field service applications. Laboratory grade oscilloscopes, especially older units that use vacuum tubes, are generally bench-top devices or are mounted on dedicated carts. Special-purpose oscilloscopes may be rack-mounted or permanently mounted into a custom instrument housing.
Inputs
The signal to be measured is fed to one of the input connectors, which is usually a coaxial connector such as a BNC or UHF type. Binding posts or banana plugs may be used for lower frequencies. If the signal source has its own coaxial connector, then a simple coaxial cable is used; otherwise, a specialized cable called a "scope probe", supplied with the oscilloscope, is used. In general, for routine use, an open wire test lead for connecting to the point being observed is not satisfactory, and a probe is generally necessary. General-purpose oscilloscopes usually present an input impedance of 1 megohm in parallel with a small but known capacitance such as 20 picofarads. This allows the use of standard oscilloscope probes. Scopes for use with very high frequencies may have 50 Ω inputs. These must be either connected directly to a 50 Ω signal source or used with Z0 or active probes.
Less-frequently-used inputs include one (or two) for triggering the sweep, horizontal deflection for X‑Y mode displays, and trace brightening/darkening, sometimes called z'‑axis inputs.
Probes
Open wire test leads (flying leads) are likely to pick up interference, so they are not suitable for low level signals. Furthermore, the leads have a high inductance, so they are not suitable for high frequencies. Using a shielded cable (i.e., coaxial cable) is better for low level signals. Coaxial cable also has lower inductance, but it has higher capacitance: a typical 50 ohm cable has about 90 pF per meter. Consequently, a one-meter direct (1×) coaxial probe loads a circuit with a capacitance of about 110 pF and a resistance of 1 megohm.
To minimize loading, attenuator probes (e.g., 10× probes) are used. A typical probe uses a 9 megohm series resistor shunted by a low-value capacitor to make an RC compensated divider with the cable capacitance and scope input. The RC time constants are adjusted to match. For example, the 9 megohm series resistor is shunted by a 12.2 pF capacitor for a time constant of 110 microseconds. The cable capacitance of 90 pF in parallel with the scope input of 20 pF and 1 megohm (total capacitance 110 pF) also gives a time constant of 110 microseconds. In practice, there is an adjustment so the operator can precisely match the low frequency time constant (called compensating the probe). Matching the time constants makes the attenuation independent of frequency. At low frequencies (where the resistance of R is much less than the reactance of C), the circuit looks like a resistive divider; at high frequencies (resistance much greater than reactance), the circuit looks like a capacitive divider.
The result is a frequency compensated probe for modest frequencies. It presents a load of about 10 megohms shunted by 12 pF. Such a probe is an improvement, but does not work well when the time scale shrinks to several cable transit times or less (transit time is typically 5 ns). In that time frame, the cable looks like its characteristic impedance, and reflections from the transmission line mismatch at the scope input and the probe causes ringing. The modern scope probe uses lossy low capacitance transmission lines and sophisticated frequency shaping networks to make the 10× probe perform well at several hundred megahertz. Consequently, there are other adjustments for completing the compensation.
Probes with 10:1 attenuation are by far the most common; for large signals (and slightly-less capacitive loading), 100:1 probes may be used. There are also probes that contain switches to select 10:1 or direct (1:1) ratios, but the latter setting has significant capacitance (tens of pF) at the probe tip, because the whole cable's capacitance is then directly connected.
Most oscilloscopes provide for probe attenuation factors, displaying the effective sensitivity at the probe tip. Historically, some auto-sensing circuitry used indicator lamps behind translucent windows in the panel to illuminate different parts of the sensitivity scale. To do so, the probe connectors (modified BNCs) had an extra contact to define the probe's attenuation. (A certain value of resistor, connected to ground, "encodes" the attenuation.) Because probes wear out, and because the auto-sensing circuitry is not compatible between different oscilloscope makes, auto-sensing probe scaling is not foolproof. Likewise, manually setting the probe attenuation is prone to user error. Setting the probe scaling incorrectly is a common error, and throws the reading off by a factor of 10.
Special high voltage probes form compensated attenuators with the oscilloscope input. These have a large probe body, and some require partly filling a canister surrounding the series resistor with volatile liquid fluorocarbon to displace air. The oscilloscope end has a box with several waveform-trimming adjustments. For safety, a barrier disc keeps the user's fingers away from the point being examined. Maximum voltage is in the low tens of kV. (Observing a high voltage ramp can create a staircase waveform with steps at different points every repetition, until the probe tip is in contact. Until then, a tiny arc charges the probe tip, and its capacitance holds the voltage (open circuit). As the voltage continues to climb, another tiny arc charges the tip further.)
There are also current probes, with cores that surround the conductor carrying current to be examined. One type has a hole for the conductor, and requires that the wire be passed through the hole for semi-permanent or permanent mounting. However, other types, used for temporary testing, have a two-part core that can be clamped around a wire. Inside the probe, a coil wound around the core provides a current into an appropriate load, and the voltage across that load is proportional to current. This type of probe only senses AC.
A more-sophisticated probe includes a magnetic flux sensor (Hall effect sensor) in the magnetic circuit. The probe connects to an amplifier, which feeds (low frequency) current into the coil to cancel the sensed field; the magnitude of the current provides the low-frequency part of the current waveform, right down to DC. The coil still picks up high frequencies. There is a combining network akin to a loudspeaker crossover.
Front panel controls
Focus control
This control adjusts CRT focus to obtain the sharpest, most-detailed trace. In practice, focus must be adjusted slightly when observing very different signals, so it must be an external control. The control varies the voltage applied to a focusing anode within the CRT. Flat-panel displays do not need this control.
Intensity control
This adjusts trace brightness. Slow traces on CRT oscilloscopes need less, and fast ones, especially if not often repeated, require more brightness. On flat panels, however, trace brightness is essentially independent of sweep speed, because the internal signal processing effectively synthesizes the display from the digitized data.
Astigmatism
This control may instead be called "shape" or "spot shape". It adjusts the voltage on the last CRT anode (immediately next to the Y deflection plates). For a circular spot, the final anode must be at the same potential as both of the Y-plates (for a centred spot the Y-plate voltages must be the same). If the anode is made more positive, the spot becomes elliptical in the X-plane as the more negative Y-plates will repel the beam. If the anode is made more negative, the spot becomes elliptical in the Y-plane as the more positive Y-plates will attract the beam. This control may be absent from simpler oscilloscope designs or may even be an internal control. It is not necessary with flat panel displays.
Beam finder
Modern oscilloscopes have direct-coupled deflection amplifiers, which means the trace could be deflected off-screen. They also might have their beam blanked without the operator knowing it. To help in restoring a visible display, the beam finder circuit overrides any blanking and limits the beam deflection to the visible portion of the screen. Beam-finder circuits often distort the trace while activated.
Graticule
The graticule is a grid of lines that serve as reference marks for measuring the displayed trace. These markings, whether located directly on the screen or on a removable plastic filter, usually consist of a 1 cm grid with closer tick marks (often at 2 mm) on the centre vertical and horizontal axis. One expects to see ten major divisions across the screen; the number of vertical major divisions varies. Comparing the grid markings with the waveform permits one to measure both voltage (vertical axis) and time (horizontal axis). Frequency can also be determined by measuring the waveform period and calculating its reciprocal.
On old and lower-cost CRT oscilloscopes the graticule is a sheet of plastic, often with light-diffusing markings and concealed lamps at the edge of the graticule. The lamps had a brightness control. Higher-cost instruments have the graticule marked on the inside face of the CRT, to eliminate parallax errors; better ones also had adjustable edge illumination with diffusing markings. (Diffusing markings appear bright.) Digital oscilloscopes, however, generate the graticule markings on the display in the same way as the trace.
External graticules also protect the glass face of the CRT from accidental impact. Some CRT oscilloscopes with internal graticules have an unmarked tinted sheet plastic light filter to enhance trace contrast; this also serves to protect the faceplate of the CRT.
Accuracy and resolution of measurements using a graticule is relatively limited; better instruments sometimes have movable bright markers on the trace. These permit internal circuits to make more refined measurements.
Both calibrated vertical sensitivity and calibrated horizontal time are set in steps. This leads, however, to some awkward interpretations of minor divisions.
Digital oscilloscopes generate the graticule digitally. The scale, spacing, etc., of the graticule can therefore be varied, and accuracy of readings may be improved.
Timebase controls
These select the horizontal speed of the CRT's spot as it creates the trace; this process is commonly referred to as the sweep. In all but the least-costly modern oscilloscopes, the sweep speed is selectable and calibrated in units of time per major graticule division. Quite a wide range of sweep speeds is generally provided, from seconds to as fast as picoseconds (in the fastest) per division. Usually, a continuously-variable control (often a knob in front of the calibrated selector knob) offers uncalibrated speeds, typically slower than calibrated. This control provides a range somewhat greater than the calibrated steps, making any speed between the steps available.
Holdoff control
Some higher-end analog oscilloscopes have a holdoff control. This sets a time after a trigger during which the sweep circuit cannot be triggered again. It helps provide a stable display of repetitive events in which some triggers would create confusing displays. It is usually set to minimum, because a longer time decreases the number of sweeps per second, resulting in a dimmer trace. See Holdoff for a more detailed description.
Vertical sensitivity, coupling, and polarity controls
To accommodate a wide range of input amplitudes, a switch selects calibrated sensitivity of the vertical deflection. Another control, often in front of the calibrated selector knob, offers a continuously variable sensitivity over a limited range from calibrated to less-sensitive settings.
Often the observed signal is offset by a steady component, and only the changes are of interest. An input coupling switch in the "AC" position connects a capacitor in series with the input that blocks low-frequency signals and DC. However, when the signal has a fixed offset of interest, or changes slowly, the user will usually prefer "DC" coupling, which bypasses any such capacitor. Most oscilloscopes offer the DC input option. For convenience, to see where zero volts input currently shows on the screen, many oscilloscopes have a third switch position (usually labeled "GND" for ground) that disconnects the input and grounds it. Often, in this case, the user centers the trace with the vertical position control.
Better oscilloscopes have a polarity selector. Normally, a positive input moves the trace upward; the polarity selector offers an "inverting" option, in which a positive-going signal deflects the trace downward.
Vertical position control
The vertical position control moves the whole displayed trace up and down. It is used to set the no-input trace exactly on the center line of the graticule, but also permits offsetting vertically by a limited amount. With direct coupling, adjustment of this control can compensate for a limited DC component of an input.
Horizontal sensitivity control
This control is found only on more elaborate oscilloscopes; it offers adjustable sensitivity for external horizontal inputs. It is only active when the instrument is in X-Y mode, i.e. the internal horizontal sweep is turned off.
Horizontal position control
The horizontal position control moves the display sidewise. It usually sets the left end of the trace at the left edge of the graticule, but it can displace the whole trace when desired. This control also moves the X-Y mode traces sidewise in some instruments, and can compensate for a limited DC component as for vertical position.
Dual-trace controls
Each input channel usually has its own set of sensitivity, coupling, and position controls, though some four-trace oscilloscopes have only minimal controls for their third and fourth channels.
Dual-trace oscilloscopes have a mode switch to select either channel alone, both channels, or (in some) an X‑Y display, which uses the second channel for X deflection. When both channels are displayed, the type of channel switching can be selected on some oscilloscopes; on others, the type depends upon timebase setting. If manually selectable, channel switching can be free-running (asynchronous), or between consecutive sweeps. Some Philips dual-trace analog oscilloscopes had a fast analog multiplier, and provided a display of the product of the input channels.
Multiple-trace oscilloscopes have a switch for each channel to enable or disable display of the channel's trace.
Delayed-sweep controls
These include controls for the delayed-sweep timebase, which is calibrated, and often also variable. The slowest speed is several steps faster than the slowest main sweep speed, though the fastest is generally the same. A calibrated multiturn delay time control offers wide range, high resolution delay settings; it spans the full duration of the main sweep, and its reading corresponds to graticule divisions (but with much finer precision). Its accuracy is also superior to that of the display.
A switch selects display modes: Main sweep only, with a brightened region showing when the delayed sweep is advancing, delayed sweep only, or (on some) a combination mode.
Good CRT oscilloscopes include a delayed-sweep intensity control, to allow for the dimmer trace of a much-faster delayed sweep which nevertheless occurs only once per main sweep. Such oscilloscopes also are likely to have a trace separation control for multiplexed display of both the main and delayed sweeps together.
Sweep trigger controls
A switch selects the trigger source. It can be an external input, one of the vertical channels of a dual or multiple-trace oscilloscope, or the AC line (mains) frequency. Another switch enables or disables auto trigger mode, or selects single sweep, if provided in the oscilloscope. Either a spring-return switch position or a pushbutton arms single sweeps.
A trigger level control varies the voltage required to generate a trigger, and the slope switch selects positive-going or negative-going polarity at the selected trigger level.
Basic types of sweep
Triggered sweep
To display events with unchanging or slowly (visibly) changing waveforms, but occurring at times that may not be evenly spaced, modern oscilloscopes have triggered sweeps. Compared to older, simpler oscilloscopes with continuously-running sweep oscillators, triggered-sweep oscilloscopes are markedly more versatile.
A triggered sweep starts at a selected point on the signal, providing a stable display. In this way, triggering allows the display of periodic signals such as sine waves and square waves, as well as nonperiodic signals such as single pulses, or pulses that do not recur at a fixed rate.
With triggered sweeps, the scope blanks the beam and starts to reset the sweep circuit each time the beam reaches the extreme right side of the screen. For a period of time, called holdoff, (extendable by a front-panel control on some better oscilloscopes), the sweep circuit resets completely and ignores triggers. Once holdoff expires, the next trigger starts a sweep. The trigger event is usually the input waveform reaching some user-specified threshold voltage (trigger level) in the specified direction (going positive or going negative—trigger polarity).
In some cases, variable holdoff time can be useful to make the sweep ignore interfering triggers that occur before the events to be observed. In the case of repetitive, but complex waveforms, variable holdoff can provide a stable display that could not otherwise be achieved.
HoldoffTrigger holdoff defines a certain period following a trigger during which the sweep cannot be triggered again. This makes it easier to establish a stable view of a waveform with multiple edges, which would otherwise cause additional triggers.
Example
Imagine the following repeating waveform:
The green line is the waveform, the red vertical partial line represents the location of the trigger, and the yellow line represents the trigger level. If the scope was simply set to trigger on every rising edge, this waveform would cause three triggers for each cycle:
Assuming the signal is fairly high frequency, the scope display would probably look something like this:
On an actual scope, each trigger would be the same channel, so all would be the same color.
It is desirable for the scope to trigger on only one edge per cycle, so it is necessary to set the holdoff at slightly less than the period of the waveform. This prevents triggering from occurring more than once per cycle, but still lets it trigger on the first edge of the next cycle.
Automatic sweep mode
Triggered sweeps can display a blank screen if there are no triggers. To avoid this, these sweeps include a timing circuit that generates free-running triggers so a trace is always visible. This is referred to as "auto sweep" or "automatic sweep" in the controls. Once triggers arrive, the timer stops providing pseudo-triggers. The user will usually disable automatic sweep when observing low repetition rates.
Recurrent sweeps
If the input signal is periodic, the sweep repetition rate can be adjusted to display a few cycles of the waveform. Early (tube) oscilloscopes and lowest-cost oscilloscopes have sweep oscillators that run continuously, and are uncalibrated. Such oscilloscopes are very simple, comparatively inexpensive, and were useful in radio servicing and some TV servicing. Measuring voltage or time is possible, but only with extra equipment, and is quite inconvenient. They are primarily qualitative instruments.
They have a few (widely spaced) frequency ranges, and relatively wide-range continuous frequency control within a given range. In use, the sweep frequency is set to slightly lower than some submultiple of the input frequency, to display typically at least two cycles of the input signal (so all details are visible). A very simple control feeds an adjustable amount of the vertical signal (or possibly, a related external signal) to the sweep oscillator. The signal triggers beam blanking and a sweep retrace sooner than it would occur free-running, and the display becomes stable.
Single sweeps
Some oscilloscopes offer these. The user manually arms the sweep circuit (typically by a pushbutton or equivalent). "Armed" means it is ready to respond to a trigger. Once the sweep completes, it resets, and does not sweep again until re-armed. This mode, combined with an oscilloscope camera, captures single-shot events.
Types of trigger include:
external trigger, a pulse from an external source connected to a dedicated input on the scope.
edge trigger, an edge detector that generates a pulse when the input signal crosses a specified threshold voltage in a specified direction. These are the most common types of triggers; the level control sets the threshold voltage, and the slope control selects the direction (negative or positive-going). (The first sentence of the description also applies to the inputs to some digital logic circuits; those inputs have fixed threshold and polarity response.)
video trigger, also known as TV trigger, a circuit that extracts synchronizing pulses from video formats such as PAL and NTSC and triggers the timebase on every line, a specified line, every field, or every frame. This circuit is typically found in a waveform monitor device, though some better oscilloscopes include this function.
delayed trigger, which waits a specified time after an edge trigger before starting the sweep. As described under delayed sweeps, a trigger delay circuit (typically the main sweep) extends this delay to a known and adjustable interval. In this way, the operator can examine a particular pulse in a long train of pulses.
Some recent designs of oscilloscopes include more sophisticated triggering schemes; these are described toward the end of this article.
Delayed sweeps
More sophisticated analog oscilloscopes contain a second timebase for a delayed sweep. A delayed sweep provides a very detailed look at some small selected portion of the main timebase. The main timebase serves as a controllable delay, after which the delayed timebase starts. This can start when the delay expires, or can be triggered (only) after the delay expires. Ordinarily, the delayed timebase is set for a faster sweep, sometimes much faster, such as 1000:1. At extreme ratios, jitter in the delays on consecutive main sweeps degrades the display, but delayed-sweep triggers can overcome this.
The display shows the vertical signal in one of several modes: the main timebase, or the delayed timebase only, or a combination thereof. When the delayed sweep is active, the main sweep trace brightens while the delayed sweep is advancing. In one combination mode, provided only on some oscilloscopes, the trace changes from the main sweep to the delayed sweep once the delayed sweep starts, though less of the delayed fast sweep is visible for longer delays. Another combination mode multiplexes (alternates) the main and delayed sweeps so that both appear at once; a trace separation control displaces them. DSOs can display waveforms this way, without offering a delayed timebase as such.
Dual and multiple-trace oscilloscopes
Oscilloscopes with two vertical inputs, referred to as dual-trace oscilloscopes, are extremely useful and commonplace.
Using a single-beam CRT, they multiplex the inputs, usually switching between them fast enough to display two traces apparently at once. Less common are oscilloscopes with more traces; four inputs are common among these, but a few (Kikusui, for one) offered a display of the sweep trigger signal if desired. Some multi-trace oscilloscopes use the external trigger input as an optional vertical input, and some have third and fourth channels with only minimal controls. In all cases, the inputs, when independently displayed, are time-multiplexed, but dual-trace oscilloscopes often can add their inputs to display a real-time analog sum. Inverting one channel while adding them together results in a display of the differences between them, provided neither channel is overloaded. This difference mode can provide a moderate-performance differential input.)
Switching channels can be asynchronous, i.e. free-running, with respect to the sweep frequency; or it can be done after each horizontal sweep is complete. Asynchronous switching is usually designated "Chopped", while sweep-synchronized is designated "Alt[ernate]". A given channel is alternately connected and disconnected, leading to the term "chopped". Multi-trace oscilloscopes also switch channels either in chopped or alternate modes.
In general, chopped mode is better for slower sweeps. It is possible for the internal chopping rate to be a multiple of the sweep repetition rate, creating blanks in the traces, but in practice this is rarely a problem. The gaps in one trace are overwritten by traces of the following sweep. A few oscilloscopes had a modulated chopping rate to avoid this occasional problem. Alternate mode, however, is better for faster sweeps.
True dual-beam CRT oscilloscopes did exist, but were not common. One type (Cossor, U.K.) had a beam-splitter plate in its CRT, and single-ended deflection following the splitter. Others had two complete electron guns, requiring tight control of axial (rotational) mechanical alignment in manufacturing the CRT. Beam-splitter types had horizontal deflection common to both vertical channels, but dual-gun oscilloscopes could have separate time bases, or use one time base for both channels. Multiple-gun CRTs (up to ten guns) were made in past decades. With ten guns, the envelope (bulb) was cylindrical throughout its length. (Also see "CRT Invention" in Oscilloscope history.)
The vertical amplifier
In an analog oscilloscope, the vertical amplifier acquires the signal[s] to be displayed and provides a signal large enough to deflect the CRT's beam. In better oscilloscopes, it delays the signal by a fraction of a microsecond. The maximum deflection is at least somewhat beyond the edges of the graticule, and more typically some distance off-screen. The amplifier has to have low distortion to display its input accurately (it must be linear), and it has to recover quickly from overloads. As well, its time-domain response has to represent transients accurately—minimal overshoot, rounding, and tilt of a flat pulse top.
A vertical input goes to a frequency-compensated step attenuator to reduce large signals to prevent overload. The attenuator feeds one or more low-level stages, which in turn feed gain stages (and a delay-line driver if there is a delay). Subsequent gain stages lead to the final output stage, which develops a large signal swing (tens of volts, sometimes over 100 volts) for CRT electrostatic deflection.
In dual and multiple-trace oscilloscopes, an internal electronic switch selects the relatively low-level output of one channel's early-stage amplifier and sends it to the following stages of the vertical amplifier.
In free-running ("chopped") mode, the oscillator (which may be simply a different operating mode of the switch driver) blanks the beam before switching, and unblanks it only after the switching transients have settled.
Part way through the amplifier is a feed to the sweep trigger circuits, for internal triggering from the signal. This feed would be from an individual channel's amplifier in a dual or multi-trace oscilloscope, the channel depending upon the setting of the trigger source selector.
This feed precedes the delay (if there is one), which allows the sweep circuit to unblank the CRT and start the forward sweep, so the CRT can show the triggering event. High-quality analog delays add a modest cost to an oscilloscope, and are omitted in cost-sensitive oscilloscopes.
The delay, itself, comes from a special cable with a pair of conductors wound around a flexible, magnetically soft core. The coiling provides distributed inductance, while a conductive layer close to the wires provides distributed capacitance. The combination is a wideband transmission line with considerable delay per unit length. Both ends of the delay cable require matched impedances to avoid reflections.
X-Y mode
Most modern oscilloscopes have several inputs for voltages, and thus can be used to plot one varying voltage versus another. This is especially useful for graphing I-V curves (current versus voltage characteristics) for components such as diodes, as well as Lissajous figures. Lissajous figures are an example of how an oscilloscope can be used to track phase differences between multiple input signals. This is very frequently used in broadcast engineering to plot the left and right stereophonic channels, to ensure that the stereo generator is calibrated properly. Historically, stable Lissajous figures were used to show that two sine waves had a relatively simple frequency relationship, a numerically-small ratio. They also indicated phase difference between two sine waves of the same frequency.
The X-Y mode also lets the oscilloscope serve as a vector monitor to display images or user interfaces. Many early games, such as Tennis for Two, used an oscilloscope as an output device.
Complete loss of signal in an X-Y CRT display means that the beam is stationary, striking a small spot. This risks burning the phosphor if the brightness is too high. Such damage was more common in older scopes as the phosphors previously used burned more easily. Some dedicated X-Y displays reduce beam current greatly, or blank the display entirely, if there are no inputs present.
Z input
Some analogue oscilloscopes feature a Z input. This is generally an input terminal that connects directly to the CRT grid (usually via a coupling capacitor). This allows an external signal to either increase (if positive) or decrease (if negative) the brightness of the trace, even allowing it to be totally blanked. The voltage range to achieve cut-off to a brightened display is of the order of 10–20 volts depending on the CRT characteristics.
An example of a practical application is if a pair of sine waves of known frequency are used to generate a circular Lissajous figure and a higher unknown frequency is applied to the Z input. This turns the continuous circle into a circle of dots. The number of dots multiplied by the X-Y frequency gives the Z frequency. This technique only works if the Z frequency is an integer ratio of the X-Y frequency and only if it is not so large that the dots become so numerous that they are difficult to count.
Bandwidth
As with all practical instruments, oscilloscopes do not respond equally to all possible input frequencies. The range of sinusoid frequencies an oscilloscope can usefully display is referred to as its bandwidth. Bandwidth applies primarily to the Y-axis, though the X-axis sweeps must be fast enough to show the highest-frequency waveforms.
The bandwidth is defined as the frequency at which the sensitivity is 0.707 of the sensitivity at DC or the lowest AC frequency (a drop of 3 dB). The oscilloscope's response drops off rapidly as the input frequency rises above that point. Within the stated bandwidth the response is not necessarily exactly uniform (or "flat"), but should always fall within a +0 to −3 dB range. One source says there is a noticeable effect on the accuracy of voltage measurements at only 20 percent of the stated bandwidth. Some oscilloscopes' specifications do include a narrower tolerance range within the stated bandwidth.
Probes also have bandwidth limits and must be chosen and used to handle the frequencies of interest properly. To achieve the flattest response, most probes must be "compensated" (an adjustment performed using a test signal from the oscilloscope) to allow for the reactance of the probe's cable.
Another related specification is rise time. This is the time taken between 10% and 90% of the maximum amplitude response at the leading edge of a pulse. It is related to the bandwidth approximately by:
Bandwidth in Hz × rise time in seconds = 0.35.
For example, an oscilloscope with a rise time of 1 nanosecond would have a bandwidth of 350 MHz.
In analog instruments, the bandwidth of the oscilloscope is limited by the vertical amplifiers and the CRT or other display subsystem. In digital instruments, the sampling rate of the analog-to-digital converter (ADC) is a factor, but the stated analog bandwidth (and therefore the overall bandwidth of the instrument) is usually less than the ADC's Nyquist frequency. This is due to limitations in the analog signal amplifier, deliberate design of the anti-aliasing filter that precedes the ADC, or both.
For a digital oscilloscope, a rule of thumb is that the continuous sampling rate should be ten times the highest frequency desired to resolve; for example a 20 megasample/second rate would be applicable for measuring signals up to about 2 MHz. This lets the anti-aliasing filter be designed with a 3 dB down point of 2 MHz and an effective cutoff at 10 MHz (the Nyquist frequency), avoiding the artifacts of a very steep ("brick-wall") filter.
A sampling oscilloscope can display signals of considerably higher frequency than the sampling rate if the signals are exactly, or nearly, repetitive. It does this by taking one sample from each successive repetition of the input waveform, each sample being at an increased time interval from the trigger event. The waveform is then displayed from these collected samples. This mechanism is referred to as "equivalent-time sampling". Some oscilloscopes can operate in either this mode or in the more traditional "real-time" mode at the operator's choice.
Waveform interval and sampling interval
For digital oscilloscopes, waveform interval is defined as the time interval between adjacent points of a displayed waveform while sampling interval is defined as the time interval between adjacent gathered samples (= 1 / sampling frequency), and the waveform interval is usually longer than the sample interval. In other words, the displayed waveform is an aggregation of the gathered samples (e.g., each displayed point is the average over each waveform interval).
Other features
Some oscilloscopes have cursors. These are lines that can be moved about the screen to measure the time interval between two points, or the difference between two voltages. A few older oscilloscopes simply brightened the trace at movable locations. These cursors are more accurate than visual estimates referring to graticule lines.
Better quality general purpose oscilloscopes include a calibration signal for setting up the compensation of test probes; this is (often) a 1 kHz square-wave signal of a definite peak-to-peak voltage available at a test terminal on the front panel. Some better oscilloscopes also have a squared-off loop for checking and adjusting current probes.
Sometimes a user wants to see an event that happens only occasionally. To catch these events, some oscilloscopes—called storage scopes—preserve the most recent sweep on the screen. This was originally achieved with a special CRT, a storage tube, which retained the image of even a very brief event for a long time.
Some digital oscilloscopes can sweep at speeds as slow as once per hour, emulating a strip chart recorder.
That is, the signal scrolls across the screen from right to left. Most oscilloscopes with this facility switch from a sweep to a strip-chart mode at about one sweep per ten seconds. This is because otherwise, the scope looks broken: it is collecting data, but the dot cannot be seen.
All but the simplest models of current oscilloscopes more often use digital signal sampling. Samples feed fast analog-to-digital converters, following which all signal processing (and storage) is digital.
Many oscilloscopes accommodate plug-in modules for different purposes, e.g., high-sensitivity amplifiers of relatively narrow bandwidth, differential amplifiers, amplifiers with four or more channels, sampling plugins for repetitive signals of very high frequency, and special-purpose plugins, including audio/ultrasonic spectrum analyzers, and stable-offset-voltage direct-coupled channels with relatively high gain.
Examples of use
One of the most frequent uses of scopes is troubleshooting malfunctioning electronic equipment. For example, where a voltmeter may show a totally unexpected voltage, a scope may reveal that the circuit is oscillating. In other cases the precise shape or timing of a pulse is important.
In a piece of electronic equipment, for example, the connections between stages (e.g., electronic mixers, electronic oscillators, amplifiers) may be 'probed' for the expected signal, using the scope as a simple signal tracer. If the expected signal is absent or incorrect, some preceding stage of the electronics is not operating correctly. Since most failures occur because of a single faulty component, each measurement can show that some of the stages of a complex piece of equipment either work, or probably did not cause the fault.
Once the faulty stage is found, further probing can usually tell a skilled technician exactly which component has failed. Once the component is replaced, the unit can be restored to service, or at least the next fault can be isolated. This sort of troubleshooting is typical of radio and TV receivers, as well as audio amplifiers, but can apply to quite different devices such as electronic motor drives.
Another use is to check newly designed circuitry. Often, a newly designed circuit misbehaves because of design errors, bad voltage levels, electrical noise etc. Digital electronics usually operate from a clock, so a dual-trace scope showing both the clock signal and a test signal dependent upon the clock is useful. Storage scopes are helpful for "capturing" rare electronic events that cause defective operation.
Oscilloscopes are often used during real-time software development to check, among other things, missed deadlines and worst-case latencies.
Pictures of use
Automotive use
First appearing in the 1970s for ignition system analysis, automotive oscilloscopes are becoming an important workshop tool for testing sensors and output signals on electronic engine management systems, braking and stability systems. Some oscilloscopes can trigger and decode serial bus messages, such as the CAN bus commonly used in automotive applications.
Software
Many oscilloscopes today provide one or more external interfaces to allow remote instrument control by external software. These interfaces (or buses) include GPIB, Ethernet, serial port, USB and Wi-Fi.
Types and models
The following section is a brief summary of various types and models available. For a detailed discussion, refer to the other article.
Cathode-ray oscilloscope (CRO)
The earliest and simplest type of oscilloscope consisted of a CRT, a vertical amplifier, a timebase, a horizontal amplifier and a power supply. These are now called "analog" scopes to distinguish them from the "digital" scopes that became common in the 1990s and later.
Analog scopes do not necessarily include a calibrated reference grid for size measurement of waves, and they may not display waves in the traditional sense of a line segment sweeping from left to right. Instead, they could be used for signal analysis by feeding a reference signal into one axis and the signal to measure into the other axis. For an oscillating reference and measurement signal, this results in a complex looping pattern referred to as a Lissajous figure. The shape of the curve can be interpreted to identify properties of the measurement signal in relation to the reference signal and is useful across a wide range of oscillation frequencies.
Dual-beam oscilloscope
The dual-beam analog oscilloscope can display two signals simultaneously. A special dual-beam CRT generates and deflects two separate beams. Multi-trace analog oscilloscopes can simulate a dual-beam display with chop and alternate sweeps—but those features do not provide simultaneous displays. (Real-time digital oscilloscopes offer the same benefits of a dual-beam oscilloscope, but they do not require a dual-beam display.) The disadvantages of the dual trace oscilloscope are that it cannot switch quickly between traces, and cannot capture two fast transient events. A dual beam oscilloscope avoids those problems.
Analog storage oscilloscope
Trace storage is an extra feature available on some analog scopes; they used direct-view storage CRTs. Storage allows a trace pattern that normally would decay in a fraction of a second to remain on the screen for several minutes or longer. An electrical circuit can then be deliberately activated to store and erase the trace on the screen.
Digital oscilloscopes
While analog devices use continually varying voltages, digital devices use numbers that correspond to samples of the voltage. In the case of digital oscilloscopes, an analog-to-digital converter (ADC) changes the measured voltages into digital information.
The digital storage oscilloscope, or DSO for short, is the standard type of oscilloscope today for the majority of industrial applications, and thanks to the low costs of entry-level oscilloscopes even for hobbyists. It replaces the electrostatic storage method in analog storage scopes with digital memory, which stores sample data as long as required without degradation and displays it without the brightness issues of storage-type CRTs. It also allows complex processing of the signal by high-speed digital signal processing circuits.
A standard DSO is limited to capturing signals with a bandwidth of less than half the sampling rate of the ADC (called the Nyquist limit). There is a variation of the DSO called the digital sampling oscilloscope which can exceed this limit for certain types of signal, such as high-speed communications signals, where the waveform consists of repeating pulses. This type of DSO deliberately samples at a much lower frequency than the Nyquist limit and then uses signal processing to reconstruct a composite view of a typical pulse.
Mixed-signal oscilloscopes
A logic analyzer is similar to an oscilloscope, but for each input signal only provides the logic level without the shape of its analog waveform. A mixed-signal oscilloscope (or MSO) meanwhile has two kinds of inputs: a small number of analog channels (typically two or four), and a larger number of logic channels (typically sixteen). It provides the ability to accurately time-correlate analog and logic signals, thus offering a distinct advantage over a separate oscilloscope and logic analyzer. Typically, logic channels may be grouped and displayed as a bus with each bus value displayed at the bottom of the display in hexadecimal or binary. On most MSOs, the trigger can be set across both analog and logic channels.
Mixed-domain oscilloscopes
A mixed-domain oscilloscope (MDO) is an oscilloscope that comes with an additional RF input which is solely used for dedicated FFT-based spectrum analyzer functionality. Often, this RF input offers a higher bandwidth than the conventional analog input channels. This is in contrast to the FFT functionality of conventional digital oscilloscopes, which use the normal analog inputs.
Some MDOs allow time-correlation of events in the time domain (like a specific serial data package) with events happening in the frequency domain (like RF transmissions).
Handheld oscilloscopes
Handheld oscilloscopes are useful for many test and field service applications. Today, a handheld oscilloscope is usually a digital sampling oscilloscope, using a liquid crystal display.
Many handheld and bench oscilloscopes have the ground reference voltage common to all input channels. If more than one measurement channel is used at the same time, all the input signals must have the same voltage reference, and the shared default reference is the "earth". If there is no differential preamplifier or external signal isolator, this traditional desktop oscilloscope is not suitable for floating measurements. (Occasionally an oscilloscope user breaks the ground pin in the power supply cord of a bench-top oscilloscope in an attempt to isolate the signal common from the earth ground. This practice is unreliable since the entire stray capacitance of the instrument cabinet connects into the circuit. It is also a hazard to break a safety ground connection, and instruction manuals strongly advise against it.)
Some models of oscilloscope have isolated inputs, where the signal reference level terminals are not connected together. Each input channel can be used to make a "floating" measurement with an independent signal reference level. Measurements can be made without tying one side of the oscilloscope input to the circuit signal common or ground reference.
The isolation available is categorized as shown below:
PC-based oscilloscopes
Some digital oscilloscope rely on a PC platform for display and control of the instrument. This can be in the form of a standalone oscilloscope with internal PC platform (PC mainboard), or as external oscilloscope which connects through USB or LAN to a separate PC or laptop.
Related instruments
A large number of instruments used in a variety of technical fields are really oscilloscopes with inputs, calibration, controls, display calibration, etc., specialized and optimized for a particular application. Examples of such oscilloscope-based instruments include waveform monitors for analyzing video levels in television productions and medical devices such as vital function monitors and electrocardiogram and electroencephalogram instruments. In automobile repair, an ignition analyzer is used to show the spark waveforms for each cylinder. All of these are essentially oscilloscopes, performing the basic task of showing the changes in one or more input signals over time in an X‑Y display.
Other instruments convert the results of their measurements to a repetitive electrical signal, and incorporate an oscilloscope as a display element. Such complex measurement systems include spectrum analyzers, transistor analyzers, and time domain reflectometers (TDRs). Unlike an oscilloscope, these instruments automatically generate stimulus or sweep a measurement parameter.
| Technology | Basics_4 | null |
14233508 | https://en.wikipedia.org/wiki/Neuston | Neuston | Neuston, also called pleuston, are organisms that live at the surface of a body of water, such as an ocean, estuary, lake, river, wetland or pond. Neuston can live on top of the water surface or submersed just below the water surface. In addition, microorganisms can exist in the surface microlayer that forms between the top- and the under-side of the water surface. Neuston has been defined as "organisms living at the air/water interface of freshwater, estuarine, and marine habitats or referring to the biota on or directly below the water’s surface layer."
Neustons can be informally separated into two groups: the phytoneuston, which are autotrophs floating at the water surface including cyanobacteria, filamentous algae and free-floating aquatic plant (e.g. mosquito fern, duckweed and water lettuce); and the zooneuston, which are floating heterotrophs such as protists (e.g. ciliates) and metazoans (aquatic animals). This article mainly concerns with metazoan zooneustons.
The word "neuston" comes from Greek neustos, meaning "swimming", and the noun suffix -on (as in "plankton"). This term first appears in the biological literature in 1917. The alternative term pleuston comes from the Greek plein, meaning "to sail or float". The first known use of this word was in 1909, before the first known use of neuston. In the past various authors have attempted distinctions between neuston and pleuston, but these distinctions have not been widely adopted. As of 2021, the two terms are usually used somewhat interchangeably, and neuston is used more often than pleuston.
Overview
The neuston of the surface layer is one of the lesser known aquatic ecological groups. The term was first used in 1917 by Naumann to describe species associated with the surface layer of freshwater habitats. Later in 1971, Zaitsev identified neuston composition in marine waters. These populations would include microscopic species, plus various plant and animal taxa, such as phytoplankton and zooplankton, living in this region. In 2002, Gladyshev further characterised the major physical and chemical dynamics of the surface layer influencing the composition and relationships with various neustonic populations"
The neustonic community structure is conditioned by sunlight and an array of endogenous (organic matter, respiratory, photosynthetic, decompositional processes) and exogenous (atmospheric deposition, inorganic matter, winds, wave action, precipitation, UV radiation, oceanic currents, surface temperature) variables and processes affecting nutrient inputs and recycling. Furthermore, the neuston provides a food source to the zooplankton migrating from deeper layers to the surface, as well as to seabirds roaming over the oceans. For these reasons, the neustonic community is believed to play a critical role on the structure and function of marine food webs. Yet, research on neuston communities to date focused predominantly on geographically limited regions of the ocean or coastal areas. Consequently, neuston complexity is still poorly understood as studies on the community structure and the taxonomical composition of organisms inhabiting this ecological niche remain few, and global scale analyses are yet lacking.
Types
There are different ways neuston can be categorised. Kennish divides them by their physical position into two groups:
epineuston: organisms living on the water's surface
hyponeuston: organisms within a region of specified depth directly below the surface layer
To this can be added the organisms living in the microlayer at the interface between air and water:
microlayer neuston: organisms (microorganisms) living in the surface microlayer sandwiched between the upper and under surface.
Marshall and Burchardt divide neuston into three ecological categories:
euneuston: organisms with maximum abundance in the vicinity of the surface on which they reside day and night
facultative neuston: organisms concentrating at the surface only during certain hours of the day, usually during darkness
pseudoneuston: organisms with maximum concentrations at deeper layers but reaching the surface layer at least during certain hours.
Freshwater neuston
Freshwater neuston, organisms living at lake or pond surfaces or slow moving parts of rivers and streams, include beetles (see whirligig beetle), protozoans, bacteria and spiders (see fishing spider and diving bell spider). Springtails in the genera Podura and Sminthurides are almost exclusively neustonic, while Hypogastrura species often aggregate on pond surfaces. Water striders such as Gerris are common examples of insects that support their weight on water's surface tension.
Floods
There are different terrestrial environmental factors such as flood pulses and droughts, and these environmental factors affect species such as neuston, whether the effects lead to more or less variations in the species. When flood pulses (an abiotic factor) occur, connectivity between different aquatic environments occur. Species that live in environments with irregular flood patterns tend to have more variations, or even decrease species and variations; similar idea to what happens when droughts occur.
Red fire ants have adapted to contend with both flooding and drought conditions. If the ants sense increased water levels in their nests, they link together and form a ball or raft that floats, with the workers on the outside and the queen inside. The brood is transported to the highest surface. They are also used as the founding structure of the raft, except for the eggs and smaller larvae. Before submerging, the ants will tip themselves into the water and sever connections with the dry land. In some cases, workers may deliberately remove all males from the raft, resulting in the males drowning.
The longevity of a raft can be as long as 12 days. Ants that are trapped underwater escape by lifting themselves to the surface using bubbles which are collected from submerged substrate. Owing to their greater vulnerability to predators, red imported fire ants are significantly more aggressive when rafting. Workers tend to deliver higher doses of venom, which reduces the threat of other animals attacking. Due to this, and because a higher workforce of ants is available, rafts are potentially dangerous to those that encounter them.
Marine neuston
The marine neuston, organisms living at the ocean surface, are one of the least studied planktonic groups. Neuston occupies a restricted ecological niche and is affected by a wide range of endogenous and exogenous processes while also being a food source to zooplankton and fish migrating from the deep layers and seabirds.
The neustonic animals form a subset of the zooplankton community, which plays a pivotal role in the functioning of marine ecosystems. Zooplankton are partially responsible for the active energy flux between superficial and deep layers of the ocean. Zooplankton species composition, biomass, and secondary production influence a wide range of trophic levels in marine communities, as they constitute a link between primary production and secondary consumers. Copepods constitute the most abundant zooplankton taxon in terms of biomass and diversity worldwide. Consequently, changes in their community composition can impact the biogeochemical cycles and might be indicative of climate variability impacts on ecosystem functioning.
Historically, zooplankton assemblages research has focused mainly on taxonomic studies and those related to community structure. However, recently, research has veered toward an alternative trait-based approach, providing a perspective more focused on groups of species with analogous functional traits. This allows individuals to be classified into types characterized by the presence/absence of certain alleles of a gene, into size classes, ecological guilds, or functional groups (FGs). Functional traits are phenotypes affecting organism fitness, growth, survival, and reproductive ability. These are regulated by the expression of genes within species, and the expression of traits regulate, in turn, the species fitness under contrasting biotic and abiotic circumstances. Moreover, a specific functional trait can also develop from the interactions between other traits and environmental conditions, leading to a given trait grouping being favoured under certain conditions. Zooplankton traits can be classified in accordance to ecological functions – feeding, growth, reproduction, survival, and other characteristics such as morphology, physiology, behaviour, or life history. Particularly, feeding strategies and trophic groups are relevant to establish feeding efficiency and associated predation risk. Additionally, they facilitate the understanding of ecosystem services associated with zooplankton, such as the distribution of fisheries or biogeochemical cycling while also allowing the positioning of zooplankton taxa in the food web.
Coral-treaders are a genus of quite rare wingless marine bugs known only from coral reefs in the Indo-Pacific region. During low tide they move over water surfaces around coral atolls and reefs similar to the more familiar water-striders, staying submerged in reef crevices during high tide.
| Physical sciences | Water: General | Earth science |
3383505 | https://en.wikipedia.org/wiki/Bi-elliptic%20transfer | Bi-elliptic transfer | In astronautics and aerospace engineering, the bi-elliptic transfer is an orbital maneuver that moves a spacecraft from one orbit to another and may, in certain situations, require less delta-v than a Hohmann transfer maneuver.
The bi-elliptic transfer consists of two half-elliptic orbits. From the initial orbit, a first burn expends delta-v to boost the spacecraft into the first transfer orbit with an apoapsis at some point away from the central body. At this point a second burn sends the spacecraft into the second elliptical orbit with periapsis at the radius of the final desired orbit, where a third burn is performed, injecting the spacecraft into the desired orbit.
While they require one more engine burn than a Hohmann transfer and generally require a greater travel time, some bi-elliptic transfers require a lower amount of total delta-v than a Hohmann transfer when the ratio of final to initial semi-major axis is 11.94 or greater, depending on the intermediate semi-major axis chosen.
The idea of the bi-elliptical transfer trajectory was first published by Ary Sternfeld in 1934.
Calculation
Delta-v
The three required changes in velocity can be obtained directly from the vis-viva equation
where
is the speed of an orbiting body,
is the standard gravitational parameter of the primary body,
is the distance of the orbiting body from the primary, i.e., the radius,
is the semi-major axis of the body's orbit.
In what follows,
is the radius of the initial circular orbit,
is the radius of the final circular orbit,
is the common apoapsis radius of the two transfer ellipses and is a free parameter of the maneuver,
and are the semimajor axes of the two elliptical transfer orbits, which are given by and
Starting from the initial circular orbit with radius (dark blue circle in the figure to the right), a prograde burn (mark 1 in the figure) puts the spacecraft on the first elliptical transfer orbit (aqua half-ellipse). The magnitude of the required delta-v for this burn is
When the apoapsis of the first transfer ellipse is reached at a distance from the primary, a second prograde burn (mark 2) raises the periapsis to match the radius of the target circular orbit, putting the spacecraft on a second elliptic trajectory (orange half-ellipse). The magnitude of the required delta-v for the second burn is
Lastly, when the final circular orbit with radius is reached, a retrograde burn (mark 3) circularizes the trajectory into the final target orbit (red circle). The final retrograde burn requires a delta-v of magnitude
If , then the maneuver reduces to a Hohmann transfer (in that case can be verified to become zero). Thus the bi-elliptic transfer constitutes a more general class of orbital transfers, of which the Hohmann transfer is a special two-impulse case.
The maximal possible savings can be computed by assuming that , in which case the total simplifies to . In this case, one also speaks of a bi-parabolic transfer because the two transfer trajectories are no longer ellipses but parabolas. The transfer time increases to infinity too.
Transfer time
Like the Hohmann transfer, both transfer orbits used in the bi-elliptic transfer constitute exactly one half of an elliptic orbit. This means that the time required to execute each phase of the transfer is half the orbital period of each transfer ellipse.
Using the equation for the orbital period and the notation from above,
The total transfer time is the sum of the times required for each half-orbit. Therefore:
and finally:
Comparison with the Hohmann transfer
Delta-v
The figure shows the total required to transfer from a circular orbit of radius to another circular orbit of radius . The is shown normalized to the orbital speed in the initial orbit, , and is plotted as a function of the ratio of the radii of the final and initial orbits, ; this is done so that the comparison is general (i.e. not dependent of the specific values of and , only on their ratio).
The thick black curve indicates the for the Hohmann transfer, while the thinner colored curves correspond to bi-elliptic transfers with varying values of the parameter , defined as the apoapsis radius of the elliptic auxiliary orbit normalized to the radius of the initial orbit, and indicated next to the curves. The inset shows a close-up of the region where the bi-elliptic curves cross the Hohmann curve for the first time.
One sees that the Hohmann transfer is always more efficient if the ratio of radii is smaller than 11.94. On the other hand, if the radius of the final orbit is more than 15.58 times larger than the radius of the initial orbit, then any bi-elliptic transfer, regardless of its apoapsis radius (as long as it's larger than the radius of the final orbit), requires less than a Hohmann transfer. Between the ratios of 11.94 and 15.58, which transfer is best depends on the apoapsis distance . For any given in this range, there is a value of above which the bi-elliptic transfer is superior and below which the Hohmann transfer is better. The following table lists the value of that results in the bi-elliptic transfer being better for some selected cases.
Transfer time
The long transfer time of the bi-elliptic transfer,
is a major drawback for this maneuver. It even becomes infinite for the bi-parabolic transfer limiting case.
The Hohmann transfer takes less than half of the time because there is just one transfer half-ellipse. To be precise,
Versatility in combination maneuvers
While a bi-elliptic transfer has a small parameter window where it's strictly superior to a Hohmann Transfer in terms of delta V for a planar transfer between circular orbits, the savings is fairly small, and a bi-elliptic transfer is a far greater aid when used in combination with certain other maneuvers.
At apoapsis, the spacecraft is travelling at low orbital velocity, and significant changes in periapsis can be achieved for small delta V cost. Transfers that resemble a bi-elliptic but which incorporate a plane-change maneuver at apoapsis can dramatically save delta-V on missions where the plane needs to be adjusted as well as the altitude, versus making the plane change in low circular orbit on top of a Hohmann transfer.
Likewise, dropping periapsis all the way into the atmosphere of a planetary body for aerobraking is inexpensive in velocity at apoapsis, but permits the use of "free" drag to aid in the final circularization burn to drop apoapsis; though it adds an extra mission stage of periapsis-raising back out of the atmosphere, this may, under some parameters, cost significantly less delta V than simply dropping periapsis in one burn from circular orbit.
Example
To transfer from a circular low Earth orbit with to a new circular orbit with using a Hohmann transfer orbit requires a Δv of . However, because , it is possible to do better with a bi-elliptic transfer. If the spaceship first accelerated 3061.04 m/s, thus achieving an elliptic orbit with apogee at , then at apogee accelerated another 608.825 m/s to a new orbit with perigee at , and finally at perigee of this second transfer orbit decelerated by 447.662 m/s, entering the final circular orbit, then the total Δv would be only 4117.53 m/s, which is 16.19 m/s (0.4%) less.
The Δv saving could be further improved by increasing the intermediate apogee, at the expense of longer transfer time. For example, an apogee of (1.3 times the distance to the Moon) would result in a 1% Δv saving over a Hohmann transfer, but require a transit time of 17 days. As an impractical extreme example, an apogee of (30 times the distance to the Moon) would result in a 2% Δv saving over a Hohmann transfer, but the transfer would require 4.5 years (and, in practice, be perturbed by the gravitational effects of other Solar system bodies). For comparison, the Hohmann transfer requires 15 hours and 34 minutes.
Evidently, the bi-elliptic orbit spends more of its delta-v closer to the planet (in the first burn). This yields a higher contribution to the specific orbital energy and, due to the Oberth effect, is responsible for the net reduction in required delta-v.
| Physical sciences | Orbital mechanics | Astronomy |
3385251 | https://en.wikipedia.org/wiki/Mudrock | Mudrock | Mudrocks are a class of fine-grained siliciclastic sedimentary rocks. The varying types of mudrocks include siltstone, claystone, mudstone and shale. Most of the particles of which the stone is composed are less than and are too small to study readily in the field. At first sight, the rock types appear quite similar; however, there are important differences in composition and nomenclature.
There has been a great deal of disagreement involving the classification of mudrocks. A few important hurdles to their classification include the following:
Mudrocks are the least understood and among the most understudied sedimentary rocks to date.
Studying mudrock constituents is difficult due to their diminutive size and susceptibility to weathering on outcrops.
And most importantly, scientists accept more than one classification scheme.
Mudrocks make up 50% of the sedimentary rocks in the geologic record and are easily the most widespread deposits on Earth. Fine sediment is the most abundant product of erosion, and these sediments contribute to the overall omnipresence of mudrocks. With increased pressure over time, the platey clay minerals may become aligned, with the appearance of parallel layering (fissility). This finely bedded material that splits readily into thin layers is called shale, as distinct from mudstone. The lack of fissility or layering in mudstone may be due either to the original texture, such as the presence of non-platy grains which disrupt the fabric during consolidation, or to the disruption of layering by burrowing organisms in the sediment prior to lithification.
From the beginning of civilization, when pottery and mudbricks were made by hand, to now, mudrocks have been important. The first book on mudrocks, Geologie des Argils by Millot, was not published until 1964; however, scientists, engineers, and oil producers have understood the significance of mudrocks since the discovery of the Burgess Shale and the relatedness of mudrocks and oil. Literature on this omnipresent rock-type has been increasing in recent years, and technology continues to allow for better analysis.
Nomenclature
Mudrocks, by definition, consist of at least fifty percent mud-sized particles. Specifically, mud is composed of silt-sized particles that are between 1/16 – 1/256 ((1/16)2) of a millimeter in diameter, and clay-sized particles which are less than 1/256 millimeter.
Mudrocks contain mostly clay minerals, and quartz and feldspars. They can also contain the following particles at less than 63 micrometres: calcite, dolomite, siderite, pyrite, marcasite, heavy minerals, and even organic carbon.
There are various synonyms for fine-grained siliciclastic rocks containing fifty percent or more of its constituents less than 1/256 of a millimeter. Mudstones, shales, lutites, and argillites are common qualifiers, or umbrella terms; however, the term mudrock has increasingly become the terminology of choice by sedimentary geologists and authors.
The term "mudrock" allows for further subdivisions of siltstone, claystone, mudstone, argilite and shale. For example, a siltstone would be made of more than 50-percent grains that equate to 1/16 - 1/256 of a millimeter. "Shale" denotes fissility, which implies an ability to part easily or break parallel to stratification. Siltstone, mudstone, and claystone implies lithified, or hardened, detritus without fissility.
Overall, "mudrocks" may be the most useful qualifying term, because it allows for rocks to be divided by its greatest portion of contributing grains and their respective grain size, whether silt, clay, or mud.
Claystone
A claystone is a lithified and non-cleavable mudrock. In order for a rock to be considered a claystone, it must consist of at least fifty percent clay (phyllosilicates), whose particle measures less than 1/256 of a millimeter in size. Clay minerals are integral to mudrocks, and represent the first or second most abundant constituent by volume. They make muds cohesive and plastic, or able to flow. Clay minerals are usually very finely grained and represent the smallest particles recognized in mudrocks. However, quartz, feldspar, iron oxides, and carbonates can also weather to the sizes of typical clay mineral grains.
For a size comparison, a clay-sized particle is 1/1000 the size of a sand grain. This means a clay particle will travel 1000 times further at constant water velocity, thus requiring quieter conditions for settlement.
The formation of clay is well understood, and can come from soil, volcanic ash, and glaciation. Ancient mudrocks are another source, because they weather and disintegrate easily. Feldspar, amphiboles, pyroxenes, and volcanic glass are the principle donors of clay minerals.
Mudstone
A mudstone is a siliciclastic sedimentary rock that contains a mixture of silt- and clay-sized particles (at least 1/3 of each).
The terminology of "mudstone" is not to be confused with the Dunham classification scheme for limestones. In Dunham's classification, a mudstone is any limestone containing less than ten percent carbonate grains. Note, a siliciclastic mudstone does not deal with carbonate grains. Friedman, Sanders, and Kopaska-Merkel (1992) suggest the use of "lime mudstone" to avoid confusion with siliciclastic rocks.
Siltstone
A siltstone is a lithified, non-cleavable mudrock. In order for a rock to be named a siltstone, it must contain over fifty percent silt-sized material. Silt is any particle smaller than sand, 1/16 of a millimeter, and larger than clay, 1/256 of millimeter. Silt is believed to be the product of physical weathering, which can involve freezing and thawing, thermal expansion, and release of pressure. Physical weathering does not involve any chemical changes in the rock, and it may be best summarised as the physical breaking apart of a rock.
One of the highest proportions of silt found on Earth is in the Himalayas, where phyllites are exposed to rainfall of up to five to ten meters (16 to 33 feet) a year. Quartz and feldspar are the biggest contributors to the silt realm, and silt tends to be non-cohesive, non-plastic, but can liquefy easily.
There is a simple test that can be done in the field to determine whether a rock is a siltstone or not, and that is to put the rock to one's teeth. If the rock feels "gritty" against one's teeth, then it is a siltstone.
Shale
Shale is a fine grained, hard, laminated mudrock, consisting of clay minerals, and quartz and feldspar silt. Shale is lithified and cleavable. It must have at least 50-percent of its particles measure less than 0.062 mm. This term is confined to argillaceous, or clay-bearing, rock.
There are many varieties of shale, including calcareous and organic-rich; however, black shale, or organic-rich shale, deserves further evaluation. In order for a shale to be a black shale, it must contain more than one percent organic carbon. A good source rock for hydrocarbons can contain up to twenty percent organic carbon. Generally, black shale receives its influx of carbon from algae, which decays and forms an ooze known as sapropel. When this ooze is cooked at desired pressure, three to six kilometers (1.8 - 3.7 miles) depth, and temperature, , it will form kerogen. Kerogen can be heated, and yield up to of natural oil and gas product per ton of rock.
Slate
Slate is a hard mudstone that has undergone metamorphism, and has well-developed cleavage. It has gone through metamorphism at temperatures between , or extreme deformation. Since slate is formed in the lower realm of metamorphism, based on pressure and temperature, slate retains its stratification and can be defined as a hard, fine-grained rock.
Slate is often used for roofing, flooring, or old-fashioned stone walls. It has an attractive appearance, and its ideal cleavage and smooth texture are desirable.
Creation of mud and mudrocks
Most mudrocks form in oceans or lakes, because these environments provide the quiet waters necessary for deposition. Although mudrocks can be found in every depositional environment on Earth, the majority are found in lakes and oceans.
Mud transport and supply
Heavy rainfall provides the kinetic motion necessary for mud, clay, and silt transport. Southeast Asia, including Bangladesh and India, receives high amounts of rain from monsoons, which then washes sediment from the Himalayas and surrounding areas to the Indian Ocean.
Warm, wet climates are best for weathering rocks, and there is more mud on ocean shelves off tropical coasts than on temperate or polar shelves. The Amazon system, for example, has the third largest sediment load on Earth, with rainfall providing clay, silt, and mud from the Andes in Peru, Ecuador, and Bolivia.
Rivers, waves, and longshore currents segregate mud, silt, and clay from sand and gravel due to fall velocity. Longer rivers, with low gradients and large watersheds, have the best carrying capacity for mud. The Mississippi River, a good example of long, low gradient river with a large amount of water, will carry mud from its northernmost sections, and deposit the material in its mud-dominated delta.
Mudrock depositional environments
Below is a listing of various environments that act as sources, modes of transportation to the oceans, and environments of deposition for mudrocks.
Alluvial environments
The Ganges in India, the Yellow in China, and the Lower Mississippi in the United States are good examples of alluvial valleys. These systems have a continuous source of water, and can contribute mud through overbank sedimentation, when mud and silt is deposited overbank during flooding, and oxbow sedimentation where an abandoned stream is filled by mud.
In order for an alluvial valley to exist there must be a highly elevated zone, usually uplifted by active tectonic movement, and a lower zone, which acts as a conduit for water and sediment to the ocean.
Glaciers
Vast quantities of mud and till are generated by glaciations and deposited on land as till and in lakes. Glaciers can erode already susceptible mudrock formations, and this process enhances glacial production of clay and silt.
The Northern Hemisphere contains 90-percent of the world's lakes larger than , and glaciers created many of those lakes. Lake deposits formed by glaciation, including deep glacial scouring, are abundant.
Non-glacial lakes
Although glaciers formed 90-percent of lakes in the Northern Hemisphere, they are not responsible for the formation of ancient lakes. Ancient lakes are the largest and deepest in the world, and hold up to twenty percent of today's petroleum reservoirs. They are also the second most abundant source of mudrocks, behind marine mudrocks.
Ancient lakes owe their abundance of mudrocks to their long lives and thick deposits. These deposits were susceptible to changes in oxygen and rainfall, and offer a robust account of paleoclimate consistency.
Deltas
A delta is a subaerial or subaqueous deposit formed where rivers or streams deposit sediment into a water body. Deltas, such as the Mississippi and Congo, have massive potential for sediment deposit, and can move sediments into deep ocean waters. Delta environments are found at the mouth of a river, where its waters slow as they enter the ocean, and silt and clay are deposited.
Low energy deltas, which deposit a great deal of mud, are located in lakes, gulfs, seas, and small oceans, where coastal currents are also low. Sand and gravel-rich deltas are high-energy deltas, where waves dominate, and mud and silt are carried much farther from the mouth of the river.
Coastlines
Coastal currents, mud supply, and waves are a key factor in coastline mud deposition. The Amazon River supplies 500 million tons of sediment, which is mostly clay, to the coastal region of northeastern South America. 250 tons of this sediment moves along the coast and is deposited. Much of the mud accumulated here is more than 20 meters (65 feet) thick, and extends into the ocean.
Much of the sediment carried by the Amazon can come from the Andes mountains, and the final distance traveled by the sediment is .
Marine environments
70-percent of the Earth's surface is covered by ocean, and marine environments are where we find the world's highest proportion of mudrocks. There is a great deal of lateral continuity found in the ocean, as opposed to continents which are confined.
In comparison, continents are temporary stewards of mud and silt, and the inevitable home of mudrock sediments is the oceans. Reference the mudrock cycle below in order to understand the burial and resurgence of the various particles.
There are various environments in the oceans, including deep-sea trenches, abyssal plains, volcanic seamounts, convergent, divergent, and transform plate margins. Not only is land a major source of the ocean sediments, but organisms living within the ocean contribute, as well.
The world's rivers transport the largest volume of suspended and dissolved loads of clay and silt to the sea, where they are deposited on ocean shelves. At the poles, glaciers and floating ice drop deposits directly to the sea floor. Winds can provide fine grained material from arid regions, and explosive volcanic eruptions contribute as well. All of these sources vary in the rate of their contribution.
Sediment moves to the deeper parts of the oceans by gravity, and the processes in the ocean are comparable to those on land.
Location has a large impact on the types of mudrocks found in ocean environments. For example, the Apalachicola River, which drains in the subtropics of the United States, carries up to sixty to eighty percent kaolinite mud, whereas the Mississippi carries only ten to twenty percent kaolinite.
The mudrock cycle
We can imagine the beginning of a mudrock's life as sediment at the top of a mountain, which may have been uplifted by plate tectonics or propelled into the air from a volcano. This sediment is exposed to rain, wind, and gravity which batters and breaks apart the rock by weathering. The products of weathering, including particles ranging from clay to silt, to pebbles and boulders, are transported to the basin below, where it can solidify into one if its many sedimentary mudstone types.
Eventually, the mudrock will move its way kilometers below the subsurface, where pressure and temperature cook the mudstone into a metamorphosed gneiss. The metamorphosed gneiss will make its way to the surface once again as country rock or as magma in a volcano, and the whole process will begin again.
Important properties
Color
Mudrocks form in various colors, including: red, purple, brown, yellow, green and grey, and even black. Shades of grey are most common in mudrocks, and darker colors of black come from organic carbons. Green mudrocks form in reducing conditions, where organic matter decomposes along with ferric iron. They can also be found in marine environments, where pelagic, or free-floating species, settle out of the water and decompose in the mudrock. Red mudrocks form when iron within the mudrock becomes oxidized, and depending on the intensity of red, one can determine if the rock has fully oxidized.
Fossils
Fossils are well preserved in mudrock formations, because the fine-grained rock protects the fossils from erosion, dissolution, and other processes of erosion. Fossils are particularly important for recording past environments. Paleontologists can look at a specific area and determine salinity, water depth, water temperature, water turbidity, and sedimentation rates with the aid of type and abundance of fossils in mudrock
One of the most famous mudrock formations is the Burgess Shale in Western Canada, which formed during the Cambrian. At this site, soft bodied creatures were preserved, some in whole, by the activity of mud in a sea. Solid skeletons are, generally, the only remnants of ancient life preserved; however, the Burgess Shale includes hard body parts such as bones, skeletons, teeth, and also soft body parts such as muscles, gills, and digestive systems. The Burgess Shale is one of the most significant fossil locations on Earth, preserving innumerable specimens of 500 million year old species, and its preservation is due to the protection of mudrock.
Another noteworthy formation is the Morrison Formation. This area covers 1.5 million square miles, stretching from Montana to New Mexico in the United States. It is considered one of the world's most significant dinosaur burial grounds, and its many fossils can be found in museums around the world. This site includes dinosaur fossils from a few dinosaur species, including the Allosaurus, Diplodocus, Stegosaurus, and Brontosaurus. There are also lungfish, freshwater mollusks, ferns and conifers. This deposit was formed by a humid, tropical climate with lakes, swamps, and rivers, which deposited mudrock. Inevitably, mudrock preserved countless specimens from the late Jurassic, roughly 150 million years ago.
Petroleum and natural gas
Mudrocks, especially black shale, are the source and containers of precious petroleum sources throughout the world. Since mudrocks and organic material require quiet water conditions for deposition, mudrocks are the most likely resource for petroleum. Mudrocks have low porosity, they are impermeable, and often, if the mudrock is not black shale, it remains useful as a seal to petroleum and natural gas reservoirs. In the case of petroleum found in a reservoir, the rock surrounding the petroleum is not the source rock, whereas black shale is a source rock.
Importance
As noted before, mudrocks make up fifty percent of the Earth's sedimentary geological record. They are widespread on Earth, and important for various industries.
Metamorphosed shale can hold emerald and gold, and mudrocks can host ore metals such as lead and zinc. Mudrocks are important in the preservation of petroleum and natural gas, due to their low porosity, and are commonly used by engineers to inhibit harmful fluid leakage from landfills.
Sandstones and carbonates record high-energy events in our history, and they are much easier to study. Interbedded between the high-energy events are mudrock formations that have recorded quieter, normal conditions in our Earth's history. It is the quieter, normal events of our geologic history we don't yet understand. Sandstones provide the big tectonic picture and some indications of water depth; mudrocks record oxygen content, a generally richer fossil abundance and diversity, and a much more informative geochemistry.
In recognition of mud and mudrocks' sometimes unappreciated importance to earth sciences, the Geological Society of London named 2015 as the "Year of Mud".
| Physical sciences | Sedimentary rocks | Earth science |
11647860 | https://en.wikipedia.org/wiki/Spacetime%20diagram | Spacetime diagram | A spacetime diagram is a graphical illustration of locations in space at various times, especially in the special theory of relativity. Spacetime diagrams can show the geometry underlying phenomena like time dilation and length contraction without mathematical equations.
The history of an object's location through time traces out a line or curve on a spacetime diagram, referred to as the object's world line. Each point in a spacetime diagram represents a unique position in space and time and is referred to as an event.
The most well-known class of spacetime diagrams are known as Minkowski diagrams, developed by Hermann Minkowski in 1908. Minkowski diagrams are two-dimensional graphs that depict events as happening in a universe consisting of one space dimension and one time dimension. Unlike a regular distance-time graph, the distance is displayed on the horizontal axis and time on the vertical axis. Additionally, the time and space units of measurement are chosen in such a way that an object moving at the speed of light is depicted as following a 45° angle to the diagram's axes.
Introduction to kinetic diagrams
Position versus time graphs
In the study of 1-dimensional kinematics, position vs. time graphs (called x-t graphs for short) provide a useful means to describe motion. Kinematic features besides the object's position are visible by the slope and shape of the lines. In Fig 1-1, the plotted object moves away from the origin at a positive constant velocity (1.66 m/s) for 6 seconds, halts for 5 seconds, then returns to the origin over a period of 7 seconds at a non-constant speed (but negative velocity).
At its most basic level, a spacetime diagram is merely a time vs position graph, with the directions of the axes in a usual p-t graph exchanged; that is, the vertical axis refers to temporal and the horizontal axis to spatial coordinate values. Especially when used in special relativity (SR), the temporal axes of a spacetime diagram are often scaled with the speed of light , and thus are often labeled by This changes the dimension of the addressed physical quantity from <Time> to <Length>, in accordance with the dimension associated with the spatial axis, which is frequently labeled
Standard configuration of reference frames
To ease insight into how spacetime coordinates, measured by observers in different reference frames, compare with each other, it is useful to standardize and simplify the setup. Two Galilean reference frames (i.e., conventional 3-space frames), S and S′ (pronounced "S prime"), each with observers O and O′ at rest in their respective frames, but measuring the other as moving with speeds ±v are said to be in standard configuration, when:
The x, y, z axes of frame S are oriented parallel to the respective primed axes of frame S′.
The origins of frames S and S′ coincide at time t = 0 in frame S and also at t′ = 0 in frame S′.
Frame S′ moves in the x-direction of frame S with velocity v as measured in frame S.
This spatial setting is displayed in the Fig 1-2, in which the temporal coordinates are separately annotated as quantities t and t'''.
In a further step of simplification it is often sufficient to consider just the direction of the observed motion and ignore the other two spatial components, allowing x and ct to be plotted in 2-dimensional spacetime diagrams, as introduced above.
Non-relativistic "spacetime diagrams"
The black axes labelled and on Fig 1-3 are the coordinate system of an observer, referred to as at rest, and who is positioned at . This observer's world line is identical with the time axis. Each parallel line to this axis would correspond also to an object at rest but at another position. The blue line describes an object moving with constant speed to the right, such as a moving observer.
This blue line labelled may be interpreted as the time axis for the second observer. Together with the axis, which is identical for both observers, it represents their coordinate system. Since the reference frames are in standard configuration, both observers agree on the location of the origin of their coordinate systems. The axes for the moving observer are not perpendicular to each other and the scale on their time axis is stretched. To determine the coordinates of a certain event, two lines, each parallel to one of the two axes, must be constructed passing through the event, and their intersections with the axes read off.
Determining position and time of the event A as an example in the diagram leads to the same time for both observers, as expected. Only for the position different values result, because the moving observer has approached the position of the event A since . Generally stated, all events on a line parallel to the axis happen simultaneously for both observers. There is only one universal time , modelling the existence of one common position axis. On the other hand, due to two different time axes the observers usually measure different coordinates for the same event. This graphical translation from and to and and vice versa is described mathematically by the so-called Galilean transformation.
Minkowski diagrams
Overview
The term Minkowski diagram refers to a specific form of spacetime diagram frequently used in special relativity. A Minkowski diagram is a two-dimensional graphical depiction of a portion of Minkowski space, usually where space has been curtailed to a single dimension. The units of measurement in these diagrams are taken such that the light cone at an event consists of the lines of slope plus or minus one through that event. The horizontal lines correspond to the usual notion of simultaneous events for a stationary observer at the origin.
A particular Minkowski diagram illustrates the result of a Lorentz transformation. The Lorentz transformation relates two inertial frames of reference, where an observer stationary at the event makes a change of velocity along the -axis. As shown in Fig 2-1, the new time axis of the observer forms an angle with the previous time axis, with . In the new frame of reference the simultaneous events lie parallel to a line inclined by to the previous lines of simultaneity. This is the new -axis. Both the original set of axes and the primed set of axes have the property that they are orthogonal with respect to the Minkowski inner product or relativistic dot product.
The original position on your time line (ct) is perpendicular to position A, the original position on your mutual timeline (x) where (t) is zero. This timeline where timelines come together are positioned then on the same timeline even when there are 2 different positions. The 2 positions are on the 45 degree Event line on the original position of A. Hence position A and position A’ on the Event line and (t)=0, relocate A’ back to position A.
Whatever the magnitude of , the line forms the universal bisector, as shown in Fig 2-2.
One frequently encounters Minkowski diagrams where the time units of measurement are scaled by a factor of such that one unit of equals one unit of . Such a diagram may have units of
Approximately 30 centimetres length and nanoseconds
Astronomical units and intervals of about 8 minutes and 19 seconds (499 seconds)
Light years and years
Light-second and second
With that, light paths are represented by lines parallel to the bisector between the axes.
Mathematical details
The angle between the and axes will be identical with that between the time axes and . This follows from the second postulate of special relativity, which says that the speed of light is the same for all observers, regardless of their relative motion (see below). The angle is given by
The corresponding boost from and to and and vice versa is described mathematically by the Lorentz transformation, which can be written
where is the Lorentz factor. By applying the Lorentz transformation, the spacetime axes obtained for a boosted frame will always correspond to conjugate diameters of a pair of hyperbolas.
As illustrated in Fig 2-3, the boosted and unboosted spacetime axes will in general have unequal unit lengths. If is the unit length on the axes of and respectively, the unit length on the axes of and is:
The -axis represents the worldline of a clock resting in , with representing the duration between two events happening on this worldline, also called the proper time between these events. Length upon the -axis represents the rest length or proper length of a rod resting in . The same interpretation can also be applied to distance upon the - and -axes for clocks and rods resting in .
History
Albert Einstein announced his theory of special relativity in 1905, with Hermann Minkowski providing his graphical representation in 1908.
In Minkowski's 1908 paper there were three diagrams, first to illustrate the Lorentz transformation, then the partition of the plane by the light-cone, and finally illustration of worldlines. The first diagram used a branch of the unit hyperbola to show the locus of a unit of proper time depending on velocity, thus illustrating time dilation. The second diagram showed the conjugate hyperbola to calibrate space, where a similar stretching leaves the impression of FitzGerald contraction. In 1914 Ludwik Silberstein included a diagram of "Minkowski's representation of the Lorentz transformation". This diagram included the unit hyperbola, its conjugate, and a pair of conjugate diameters. Since the 1960s a version of this more complete configuration has been referred to as The Minkowski Diagram, and used as a standard illustration of the transformation geometry of special relativity. E. T. Whittaker has pointed out that the principle of relativity is tantamount to the arbitrariness of what hyperbola radius is selected for time in the Minkowski diagram. In 1912 Gilbert N. Lewis and Edwin B. Wilson applied the methods of synthetic geometry to develop the properties of the non-Euclidean plane that has Minkowski diagrams.Synthetic Spacetime, a digest of the axioms used, and theorems proved, by Wilson and Lewis.
When Taylor and Wheeler composed Spacetime Physics (1966), they did not use the term Minkowski diagram for their spacetime geometry. Instead they included an acknowledgement of Minkowski's contribution to philosophy by the totality of his innovation of 1908.
Loedel diagrams
While a frame at rest in a Minkowski diagram has orthogonal spacetime axes, a frame moving relative to the rest frame in a Minkowski diagram has spacetime axes which form an acute angle. This asymmetry of Minkowski diagrams can be misleading, since special relativity postulates that any two inertial reference frames must be physically equivalent. The Loedel diagram is an alternative spacetime diagram that makes the symmetry of inertial references frames much more manifest.
Formulation via median frame
Several authors showed that there is a frame of reference between the resting and moving ones where their symmetry would be apparent ("median frame"). In this frame, the two other frames are moving in opposite directions with equal speed. Using such coordinates makes the units of length and time the same for both axes. If and are given between and , then these expressions are connected with the values in their median frame S0 as follows:
For instance, if between and , then by (2) they are moving in their median frame S0 with approximately each in opposite directions. On the other hand, if in S0, then by (1) the relative velocity between and in their own rest frames is . The construction of the axes of and is done in accordance with the ordinary method using with respect to the orthogonal axes of the median frame (Fig. 3–1).
However, it turns out that when drawing such a symmetric diagram, it is possible to derive the diagram's relations even without mentioning the median frame and at all. Instead, the relative velocity between and can directly be used in the following construction, providing the same result:
If is the angle between the axes of and (or between and ), and between the axes of and , it is given:
Two methods of construction are obvious from Fig. 3-2: the -axis is drawn perpendicular to the -axis, the and -axes are added at angle ; and the x′-axis is drawn at angle with respect to the -axis, the -axis is added perpendicular to the -axis and the -axis perpendicular to the -axis.
In a Minkowski diagram, lengths on the page cannot be directly compared to each other, due to warping factor between the axes' unit lengths in a Minkowski diagram. In particular, if and are the unit lengths of the rest frame axes and moving frame axes, respectively, in a Minkowski diagram, then the two unit lengths are warped relative to each other via the formula:
By contrast, in a symmetric Loedel diagram, both the and frame axes are warped by the same factor relative to the median frame and hence have identical unit lengths. This implies that, for a Loedel spacetime diagram, we can directly compare spacetime lengths between different frames as they appear on the page; no unit length scaling/conversion between frames is necessary due to the symmetric nature of the Loedel diagram.
History
Max Born (1920) drew Minkowski diagrams by placing the -axis almost perpendicular to the -axis, as well as the -axis to the -axis, in order to demonstrate length contraction and time dilation in the symmetric case of two rods and two clocks moving in opposite direction.
Dmitry Mirimanoff (1921) showed that there is always a median frame with respect to two relatively moving frames, and derived the relations between them from the Lorentz transformation. However, he did not give a graphical representation in a diagram.
Symmetric diagrams were systematically developed by Paul Gruner in collaboration with Josef Sauter in two papers in 1921. Relativistic effects such as length contraction and time dilation and some relations to covariant and contravariant vectors were demonstrated by them. (translation: An elementary geometrical representation of the transformation formulas of the special theory of relativity) Gruner extended this method in subsequent papers (1922–1924), and gave credit to Mirimanoff's treatment as well. (translation: Graphical representation of the four-dimensional space-time universe)
The construction of symmetric Minkowski diagrams was later independently rediscovered by several authors. For instance, starting in 1948, Enrique Loedel Palumbo published a series of papers in Spanish language, presenting the details of such an approach.Fisica relativista, Kapelusz Editorial, Buenos Aires, Argentina (1955). In 1955, Henri Amar also published a paper presenting such relations, and gave credit to Loedel in a subsequent paper in 1957. Some authors of textbooks use symmetric Minkowski diagrams, denoting as Loedel diagrams.
Relativistic phenomena in diagrams
Time dilation
Relativistic time dilation refers to the fact that a clock (indicating its proper time in its rest frame) that moves relative to an observer is observed to run slower. The situation is depicted in the symmetric Loedel diagrams of Fig 4-1. Note that we can compare spacetime lengths on page directly with each other, due to the symmetric nature of the Loedel diagram.
In Fig 4-2, the observer whose reference frame is given by the black axes is assumed to move from the origin O towards A. The moving clock has the reference frame given by the blue axes and moves from O to B. For the black observer, all events happening simultaneously with the event at A are located on a straight line parallel to its space axis. This line passes through A and B, so A and B are simultaneous from the reference frame of the observer with black axes. However, the clock that is moving relative to the black observer marks off time along the blue time axis. This is represented by the distance from O to B. Therefore, the observer at A with the black axes notices their clock as reading the distance from O to A while they observe the clock moving relative him or her to read the distance from O to B. Due to the distance from O to B being smaller than the distance from O to A, they conclude that the time passed on the clock moving relative to them is smaller than that passed on their own clock.
A second observer, having moved together with the clock from O to B, will argue that the black axis clock has only reached C and therefore runs slower. The reason for these apparently paradoxical statements is the different determination of the events happening synchronously at different locations. Due to the principle of relativity, the question of who is right has no answer and does not make sense.
Length contraction
Relativistic length contraction refers to the fact that a ruler (indicating its proper length in its rest frame) that moves relative to an observer is observed to contract/shorten. The situation is depicted in symmetric Loedel diagrams in Fig 4-3. Note that we can compare spacetime lengths on page directly with each other, due to the symmetric nature of the Loedel diagram.
In Fig 4-4, the observer is assumed again to move along the -axis. The world lines of the endpoints of an object moving relative to him are assumed to move along the -axis and the parallel line passing through A and B. For this observer the endpoints of the object at are O and A. For a second observer moving together with the object, so that for him the object is at rest, it has the proper length OB at . Due to . the object is contracted for the first observer.
The second observer will argue that the first observer has evaluated the endpoints of the object at O and A respectively and therefore at different times, leading to a wrong result due to his motion in the meantime. If the second observer investigates the length of another object with endpoints moving along the -axis and a parallel line passing through C and D he concludes the same way this object to be contracted from OD to OC. Each observer estimates objects moving with the other observer to be contracted. This apparently paradoxical situation is again a consequence of the relativity of simultaneity as demonstrated by the analysis via Minkowski diagram.
For all these considerations it was assumed, that both observers take into account the speed of light and their distance to all events they see in order to determine the actual times at which these events happen from their point of view.
Constancy of the speed of light
Another postulate of special relativity is the constancy of the speed of light. It says that any observer in an inertial reference frame measuring the vacuum speed of light relative to themself obtains the same value regardless of his own motion and that of the light source. This statement seems to be paradoxical, but it follows immediately from the differential equation yielding this, and the Minkowski diagram agrees. It explains also the result of the Michelson–Morley experiment which was considered to be a mystery before the theory of relativity was discovered, when photons were thought to be waves through an undetectable medium.
For world lines of photons passing the origin in different directions and holds. That means any position on such a world line corresponds with steps on - and -axes of equal absolute value. From the rule for reading off coordinates in coordinate system with tilted axes follows that the two world lines are the angle bisectors of the - and -axes. As shown in Fig 4-5, the Minkowski diagram illustrates them as being angle bisectors of the - and -axes as well. That means both observers measure the same speed for both photons.
Further coordinate systems corresponding to observers with arbitrary velocities can be added to this Minkowski diagram. For all these systems both photon world lines represent the angle bisectors of the axes. The more the relative speed approaches the speed of light the more the axes approach the corresponding angle bisector. The axis is always more flat and the time axis more steep than the photon world lines. The scales on both axes are always identical, but usually different from those of the other coordinate systems.
Speed of light and causality
Straight lines passing the origin which are steeper than both photon world lines correspond with objects moving more slowly than the speed of light. If this applies to an object, then it applies from the viewpoint of all observers, because the world lines of these photons are the angle bisectors for any inertial reference frame. Therefore, any point above the origin and between the world lines of both photons can be reached with a speed smaller than that of the light and can have a cause-and-effect relationship with the origin. This area is the absolute future, because any event there happens later compared to the event represented by the origin regardless of the observer, which is obvious graphically from the Minkowski diagram in Fig 4-6.
Following the same argument the range below the origin and between the photon world lines is the absolute past relative to the origin. Any event there belongs definitely to the past and can be the cause of an effect at the origin.
The relationship between any such pairs of event is called timelike, because they have a time distance greater than zero for all observers. A straight line connecting these two events is always the time axis of a possible observer for whom they happen at the same place. Two events which can be connected just with the speed of light are called lightlike.
In principle a further dimension of space can be added to the Minkowski diagram leading to a three-dimensional representation. In this case the ranges of future and past become cones with apexes touching each other at the origin. They are called light cones.
The speed of light as a limit
Following the same argument, all straight lines passing through the origin and which are more nearly horizontal than the photon world lines, would correspond to objects or signals moving faster than light regardless of the speed of the observer. Therefore, no event outside the light cones can be reached from the origin, even by a light-signal, nor by any object or signal moving with less than the speed of light. Such pairs of events are called spacelike because they have a finite spatial distance different from zero for all observers. On the other hand, a straight line connecting such events is always the space coordinate axis of a possible observer for whom they happen at the same time. By a slight variation of the velocity of this coordinate system in both directions it is always possible to find two inertial reference frames whose observers estimate the chronological order of these events to be different.
Given an object moving faster than light, say from O to A in Fig 4-7, then for any observer watching the object moving from O to A, another observer can be found (moving at less than the speed of light with respect to the first) for whom the object moves from A to O. The question of which observer is right has no unique answer, and therefore makes no physical sense. Any such moving object or signal would violate the principle of causality.
Also, any general technical means of sending signals faster than light would permit information to be sent into the originator's own past. In the diagram, an observer at O in the system sends a message moving faster than light to A. At A, it is received by another observer, moving so as to be in the system, who sends it back, again faster than light, arriving at B. But B is in the past relative to O. The absurdity of this process becomes obvious when both observers subsequently confirm that they received no message at all, but all messages were directed towards the other observer as can be seen graphically in the Minkowski diagram. Furthermore, if it were possible to accelerate an observer to the speed of light, their space and time axes would coincide with their angle bisector. The coordinate system would collapse, in concordance with the fact that due to time dilation, time would effectively stop passing for them.
These considerations show that the speed of light as a limit is a consequence of the properties of spacetime, and not of the properties of objects such as technologically imperfect space ships. The prohibition of faster-than-light motion, therefore, has nothing in particular to do with electromagnetic waves or light, but comes as a consequence of the structure of spacetime.
Accelerating observers
It is often, incorrectly, asserted that special relativity cannot handle accelerating particles or accelerating reference frames. In reality, accelerating particles present no difficulty at all in special relativity. On the other hand, accelerating frames do require some special treatment, However, as long as one is dealing with flat, Minkowskian spacetime, special relativity can handle the situation. It is only in the presence of gravitation that general relativity is required.
An accelerating particle's 4-vector acceleration is the derivative with respect to proper time of its 4-velocity. This is not a difficult situation to handle. Accelerating frames require that one understand the concept of a momentarily comoving reference frame (MCRF), which is to say, a frame traveling at the same instantaneous velocity of a particle at any given instant.
Consider the animation in Fig 5–1. The curved line represents the world line of a particle that undergoes continuous acceleration, including complete changes of direction in the positive and negative x''-directions. The red axes are the axes of the MCRF for each point along the particle's trajectory. The coordinates of events in the unprimed (stationary) frame can be related to their coordinates in any momentarily co-moving primed frame using the Lorentz transformations.
Fig 5-2 illustrates the changing views of spacetime along the world line of a rapidly accelerating particle. The axis (not drawn) is vertical, while the axis (not drawn) is horizontal. The dashed line is the spacetime trajectory ("world line") of the particle. The balls are placed at regular intervals of proper time along the world line. The solid diagonal lines are the light cones for the observer's current event, and they intersect at that event. The small dots are other arbitrary events in the spacetime.
The slope of the world line (deviation from being vertical) is the velocity of the particle on that section of the world line. Bends in the world line represent particle acceleration. As the particle accelerates, its view of spacetime changes. These changes in view are governed by the Lorentz transformations. Also note that:
the balls on the world line before/after future/past accelerations are more spaced out due to time dilation.
events which were simultaneous before an acceleration (horizontally spaced events) are at different times afterwards due to the relativity of simultaneity,
events pass through the light cone lines due to the progression of proper time, but not due to the change of views caused by the accelerations, and
the world line always remains within the future and past light cones of the current event.
If one imagines each event to be the flashing of a light, then the events that are within the past light cone of the observer are the events visible to the observer. The slope of the world line (deviation from being vertical) gives the velocity relative to the observer.
Case of non-inertial reference frames
The photon world lines are determined using the metric with . The light cones are deformed according to the position. In an inertial reference frame a free particle has a straight world line. In a non-inertial reference frame the world line of a free particle is curved.
Take the example of the fall of an object dropped without initial velocity from a rocket. The rocket has a uniformly accelerated motion with respect to an inertial reference frame. As can be seen from Fig 6-2 of a Minkowski diagram in a non-inertial reference frame, the object once dropped, gains speed, reaches a maximum, and then sees its speed decrease and asymptotically cancel on the horizon where its proper time freezes at . The velocity is measured by an observer at rest in the accelerated rocket.
| Physical sciences | Theory of relativity | Physics |
7906908 | https://en.wikipedia.org/wiki/Biomass%20%28energy%29 | Biomass (energy) | In the context of energy production, biomass is matter from recently living (but now dead) organisms which is used for bioenergy production. Examples include wood, wood residues, energy crops, agricultural residues including straw, and organic waste from industry and households. Wood and wood residues is the largest biomass energy source today. Wood can be used as a fuel directly or processed into pellet fuel or other forms of fuels. Other plants can also be used as fuel, for instance maize, switchgrass, miscanthus and bamboo. The main waste feedstocks are wood waste, agricultural waste, municipal solid waste, and manufacturing waste. Upgrading raw biomass to higher grade fuels can be achieved by different methods, broadly classified as thermal, chemical, or biochemical.
The climate impact of bioenergy varies considerably depending on where biomass feedstocks come from and how they are grown. For example, burning wood for energy releases carbon dioxide. Those emissions can be significantly offset if the trees that were harvested are replaced by new trees in a well-managed forest, as the new trees will remove carbon dioxide from the air as they grow. However, the farming of biomass feedstocks can reduce biodiversity, degrade soils and take land out of food production. It may also consume water for irrigation and fertilisers.
Terminology
Biomass (in the context of energy generation) is matter from recently living (but now dead) organisms which is used for bioenergy production. There are variations in how such biomass for energy is defined, e.g. only from plants, or from plants and algae, or from plants and animals. The vast majority of biomass used for bioenergy does come from plants. Bioenergy is a type of renewable energy with potential to assist with climate change mitigation.
Some people use the terms biomass and biofuel interchangeably, but it is now more common to consider biofuel to be a liquid or gaseous fuel used for transportation, as defined by government authorities in the US and EU. From that perspective, biofuel is a subset of biomass.
The European Union's Joint Research Centre defines solid biofuel as raw or processed organic matter of biological origin used for energy, such as firewood, wood chips, and wood pellets.
Types and uses
Different types of biomass are used for different purposes:
Primary biomass sources that are appropriate for heat or electricity generation but not for transport include: wood, wood residues, wood pellets, agricultural residues, organic waste.
Biomass that is processed into transport fuels can come from corn, sugar cane, and soy.
Biomass is categorized either as biomass harvested directly for energy (primary biomass), or as residues and waste: (secondary biomass).
Biomass harvested directly for energy
The main biomass types harvested directly for energy is wood, some food crops and all perennial energy crops. One third of the global forest area of 4 billion hectares is used for wood production or other commercial purposes, and forests provide 85% of all biomass used for energy globally. In the EU, forests provide 60% of all biomass used for energy, with wood residues and waste being the largest source.
Woody biomass used for energy often consists of trees and bushes harvested for traditional cooking and heating purposes, particularly in developing countries, with 25 EJ per year used globally for these purposes. This practice is highly polluting. The World Health Organization (WHO) estimates that cooking-related pollution causes 3.8 million annual deaths. The United Nations Sustainable Development Goal 7 aims for the traditional use of biomass for cooking to be phased out by 2030. Short-rotation coppices and short-rotation forests are also harvested directly for energy, providing 4 EJ of energy, and are considered sustainable. The potential for these crops and perennial energy crops to provide at least 25 EJ annually by 2050 is estimated.
Food crops harvested for energy include sugar-producing crops (such as sugarcane), starch-producing crops (such as maize), and oil-producing crops (such as rapeseed). Sugarcane is a perennial crop, while corn and rapeseed are annual crops. Sugar- and starch-producing crops are used to make bioethanol, and oil-producing crops are used to make biodiesel. The United States is the largest producer of bioethanol, while the European Union is the largest producer of biodiesel. The global production of bioethanol and biodiesel provides 2.2 and 1.5 EJ of energy per year, respectively. Biofuel made from food crops harvested for energy is also known as "first-generation" or "traditional" biofuel and has relatively low emission savings.
The IPCC estimates that between 0.32 and 1.4 billion hectares of marginal land are suitable for bioenergy worldwide.
Biomass in the form of residues and waste
Residues and waste are by-products from biological material harvested mainly for non-energy purposes. The most important by-products are wood residues, agricultural residues and municipal/industrial waste:
Wood residues are by-products from forestry operations or from the wood processing industry. Had the residues not been collected and used for bioenergy, they would have decayed (and therefore produced emissions) on the forest floor or in landfills, or been burnt (and produced emissions) at the side of the road in forests or outside wood processing facilities.
The by-products from forestry operations are called logging residues or forest residues, and consist of tree tops, branches, stumps, damaged or dying or dead trees, irregular or bent stem sections, thinnings (small trees that are cleared away in order to help the bigger trees grow large), and trees removed to reduce wildfire risk. The extraction level of logging residues differ from region to region, but there is an increasing interest in using this feedstock, since the sustainable potential is large (15 EJ annually). 68% of the total forest biomass in the EU consists of wood stems, and 32% consists of stumps, branches and tops.
The by-products from the wood processing industry are called wood processing residues and consist of cut offs, shavings, sawdust, bark, and black liquor. Wood processing residues have a total energy content of 5.5 EJ annually. Wood pellets are mainly made from wood processing residues, and have a total energy content of 0.7 EJ. Wood chips are made from a combination of feedstocks, and have a total energy content of 0.8 EJ.
The energy content in agricultural residues used for energy is approximately 2 EJ. However, agricultural residues has a large untapped potential. The energy content in the global production of agricultural residues has been estimated to 78 EJ annually, with the largest share from straw (51 EJ). Others have estimated between 18 and 82 EJ. The use of agricultural residues and waste that is both sustainable and economically feasible is expected to increase to between 37 and 66 EJ in 2030.
Municipal waste produced 1.4 EJ and industrial waste 1.1 EJ. Wood waste from cities and industry also produced 1.1 EJ. The sustainable potential for wood waste has been estimated to 2–10 EJ. IEA recommends a dramatic increase in waste utilization to 45 EJ annually in 2050.
Biomass conversion
Raw biomass can be upgraded into better and more practical fuel simply by compacting it (e.g. wood pellets), or by different conversions broadly classified as thermal, chemical, and biochemical. Biomass conversion reduces the transport costs as it is cheaper to transport high density commodities.
Thermal conversion
Thermal upgrading produces solid, liquid or gaseous fuels, with heat as the dominant conversion driver. The basic alternatives are torrefaction, pyrolysis, and gasification, these are separated principally by how far the chemical reactions involved are allowed to proceed. The advancement of the chemical reactions is mainly controlled by how much oxygen is available, and the conversion temperature.
Torrefaction is a mild form of pyrolysis where organic materials are heated to 400–600 °F (200–300 °C) in a no–to–low oxygen environment. The heating process removes (via gasification) the parts of the biomass that has the lowest energy content, while the parts with the highest energy content remain. That is, approximately 30% of the biomass is converted to gas during the torrefaction process, while 70% remains, usually in the form of compacted pellets or briquettes. This solid product is water resistant, easy to grind, non-corrosive, and it contains approximately 85% of the original biomass energy. Basically the mass part has shrunk more than the energy part, and the consequence is that the calorific value of torrefied biomass increases significantly, to the extent that it can compete with coals used for electricity generation (steam/thermal coals). The energy density of the most common steam coals today is 22–26 GJ/t. There are other less common, more experimental or proprietary thermal processes that may offer benefits, such as hydrothermal upgrading (sometimes called "wet" torrefaction.) The hydrothermal upgrade path can be used for both low and high moisture content biomass, e.g. aqueous slurries.
Pyrolysis entails heating organic materials to 800–900 °F (400–500 °C) in the near complete absence of oxygen. Biomass pyrolysis produces fuels such as bio-oil, charcoal, methane, and hydrogen. Hydrotreating is used to process bio-oil (produced by fast pyrolysis) with hydrogen under elevated temperatures and pressures in the presence of a catalyst to produce renewable diesel, renewable gasoline, and renewable jet fuel.
Gasification entails heating organic materials to 1,400–1700 °F (800–900 °C) with injections of controlled amounts of oxygen and/or steam into the vessel to produce a carbon monoxide and hydrogen rich gas called synthesis gas or syngas. Syngas can be used as a fuel for diesel engines, for heating, and for generating electricity in gas turbines. It can also be treated to separate the hydrogen from the gas, and the hydrogen can be burned or used in fuel cells. The syngas can be further processed to produce liquid fuels using the Fischer-Tropsch synthesis process.
Chemical conversion
A range of chemical processes may be used to convert biomass into other forms, such as to produce a fuel that is more practical to store, transport and use, or to exploit some property of the process itself. Many of these processes are based in large part on similar coal-based processes, such as the Fischer-Tropsch synthesis. A chemical conversion process known as transesterification is used for converting vegetable oils, animal fats, and greases into fatty acid methyl esters (FAME), which are used to produce biodiesel.
Biochemical conversion
Biochemical processes have developed in nature to break down the molecules of which biomass is composed, and many of these can be harnessed. In most cases, microorganisms are used to perform the conversion. The processes are called anaerobic digestion, fermentation, and composting.
Fermentation converts biomass into bioethanol, and anaerobic digestion converts biomass into renewable natural gas (biogas). Bioethanol is used as a vehicle fuel. Renewable natural gas—also called biogas or biomethane—is produced in anaerobic digesters at sewage treatment plants and at dairy and livestock operations. It also forms in and may be captured from solid waste landfills. Properly treated renewable natural gas has the same uses as fossil fuel natural gas.
Climate impacts
Short-term vs long-term climate benefits
Regarding the issue of climate consequences for modern bioenergy, IPCC states: "Life-cycle GHG emissions of modern bioenergy alternatives are usually lower than those for fossil fuels." Consequently, most of IPCC's GHG mitigation pathways include substantial deployment of bioenergy technologies.
Some research groups state that even if the European and North American forest carbon stock is increasing, it simply takes too long for harvested trees to grow back. Bioenergy from sources with high payback and parity times take a long time to have an impact on climate change mitigation. They therefore suggest that the EU should adjust its sustainability criteria so that only renewable energy with carbon payback times of less than 10 years is defined as sustainable, for instance wind, solar, biomass from wood residues and tree thinnings that would otherwise be burnt or decompose relatively fast, and biomass from short rotation coppicing (SRC).
The IPCC states: "While individual stands in a forest may be either sources or sinks, the forest carbon balance is determined by the sum of the net balance of all stands." IPCC also state that the only universally applicable approach to carbon accounting is the one that accounts for both carbon emissions and carbon removals (absorption) for managed lands (e.g. forest landscapes.) When the total is calculated, natural disturbances like fires and insect infestations are subtracted, and what remains is the human influence.
IEA Bioenergy state that an exclusive focus on the short-term make it harder to achieve efficient carbon mitigation in the long term, and compare investments in new bioenergy technologies with investments in other renewable energy technologies that only provide emission reductions after 2030, for instance the scaling-up of battery manufacturing or the development of rail infrastructure. Forest carbon emission avoidance strategies give a short-term mitigation benefit, but the long-term benefits from sustainable forestry activities provide ongoing forest product and energy resources.
Most of IPCC's GHG mitigation pathways include substantial deployment of bioenergy technologies. Limited or no bioenergy pathways leads to increased climate change or shifting bioenergy's mitigation load to other sectors. In addition, mitigation cost increases.
Carbon accounting system boundaries
Carbon positive scenarios are likely to be net emitters of CO2, carbon negative projects are net absorbers of CO2, while carbon neutral projects balance emissions and absorption equally.
It is common to include alternative scenarios (also called "reference scenarios" or "counterfactuals") for comparison. The alternative scenarios range from scenarios with only modest changes compared to the existing project, all the way to radically different ones (i.e. forest protection or "no-bioenergy" counterfactuals.) Generally, the difference between scenarios is seen as the actual carbon mitigation potential of the scenarios.
In addition to the choice of alternative scenario, other choices has to be made as well. The so-called "system boundaries" determine which carbon emissions/absorptions that will be included in the actual calculation, and which that will be excluded. System boundaries include temporal, spatial, efficiency-related and economic boundaries.
For example, the actual carbon intensity of bioenergy varies with biomass production techniques and transportation lengths.
Temporal system boundaries
The temporal boundaries define when to start and end carbon counting. Sometimes "early" events are included in the calculation, for instance carbon absorption going on in the forest before the initial harvest. Sometimes "late" events are included as well, for instance emissions caused by end-of-life activities for the infrastructure involved, e.g. demolition of factories. Since the emission and absorption of carbon related to a project or scenario changes with time, the net carbon emission can either be presented as time-dependent (for instance a curve which moves along a time axis), or as a static value; this shows average emissions calculated over a defined time period.
The time-dependent net emission curve will typically show high emissions at the beginning (if the counting starts when the biomass is harvested.) Alternatively, the starting point can be moved back to the planting event; in this case the curve can potentially move below zero (into carbon negative territory) if there is no carbon debt from land use change to pay back, and in addition more and more carbon is absorbed by the planted trees. The emission curve then spikes upward at harvest. The harvested carbon is then being distributed into other carbon pools, and the curve moves in tandem with the amount of carbon that is moved into these new pools (Y axis), and the time it takes for the carbon to move out of the pools and return to the forest via the atmosphere (X axis). As described above, the carbon payback time is the time it takes for the harvested carbon to be returned to the forest, and the carbon parity time is the time it takes for the carbon stored in two competing scenarios to reach the same level.
The static carbon emission value is produced by calculating the average annual net emission for a specific time period. The specific time period can be the expected lifetime of the infrastructure involved (typical for life cycle assessments; LCA's), policy relevant time horizons inspired by the Paris agreement (for instance remaining time until 2030, 2050 or 2100), time spans based on different global warming potentials (GWP; typically 20 or 100 years), or other time spans. In the EU, a time span of 20 years is used when quantifying the net carbon effects of a land use change. Generally in legislation, the static number approach is preferred over the dynamic, time-dependent curve approach. The number is expressed as a so-called "emission factor" (net emission per produced energy unit, for instance kg CO2e per GJ), or even simpler as an average greenhouse gas savings percentage for specific bioenergy pathways. The EU's published greenhouse gas savings percentages for specific bioenergy pathways used in the Renewable Energy Directive (RED) and other legal documents are based on life cycle assessments (LCA's).
Spatial system boundaries
The spatial boundaries define "geographical" borders for carbon emission/absorption calculations. The two most common spatial boundaries for CO2 absorption and emission in forests are 1.) along the edges of a particular forest stand and 2.) along the edges of a whole forest landscape, which include many forest stands of increasing age (the forest stands are harvested and replanted, one after the other, over as many years as there are stands.) A third option is the so-called increasing stand level carbon accounting method. The researcher has to decide whether to focus on the individual stand, an increasing number of stands, or the whole forest landscape. The IPCC recommends landscape-level carbon accounting.
Further, the researcher has to decide whether emissions from direct/indirect land use change should be included in the calculation. Most researchers include emissions from direct land use change, for instance the emissions caused by cutting down a forest in order to start some agricultural project there instead. The inclusion of indirect land use change effects is more controversial, as they are difficult to quantify accurately. Other choices involve defining the likely spatial boundaries of forests in the future.
Efficiency-related system boundaries
The efficiency-related boundaries define a range of fuel substitution efficiencies for different biomass-combustion pathways. Different supply chains emit different amounts of carbon per supplied energy unit, and different combustion facilities convert the chemical energy stored in different fuels to heat or electrical energy with different efficiencies. The researcher has to know about this and choose a realistic efficiency range for the different biomass-combustion paths under consideration. The chosen efficiencies are used to calculate so-called "displacement factors" – single numbers that shows how efficient fossil carbon is substituted by biogenic carbon. If for instance 10 tonnes of carbon are combusted with an efficiency half that of a modern coal plant, only 5 tonnes of coal would actually be counted as displaced (displacement factor 0.5).
Generally, fuel burned in inefficient (old or small) combustion facilities gets assigned lower displacement factors than fuel burned in efficient (new or large) facilities, since more fuel has to be burned (and therefore more CO2 released) in order to produce the same amount of energy.
The displacement factor varies with the carbon intensity of both the biomass fuel and the displaced fossil fuel. If or when bioenergy can achieve negative emissions (e.g. from afforestation, energy grass plantations and/or bioenergy with carbon capture and storage (BECCS), or if fossil fuel energy sources with higher emissions in the supply chain start to come online (e.g. because of fracking, or increased use of shale gas), the displacement factor will start to rise. On the other hand, if or when new baseload energy sources with lower emissions than fossil fuels start to come online, the displacement factor will start to drop. Whether a displacement factor change is included in the calculation or not, depends on whether or not it is expected to take place within the time period covered by the relevant scenario's temporal system boundaries.
Economic system boundaries
The economic boundaries define which market effects to include in the calculation, if any. Changed market conditions can lead to small or large changes in carbon emissions and absorptions from supply chains and forests, for instance changes in forest area as a response to changes in demand. Macroeconomic events/policy changes can have impacts on forest carbon stock. Like with indirect land use changes, economic changes can be difficult to quantify however, so some researchers prefer to leave them out of the calculation.
System boundary impacts
The chosen system boundaries are very important for the calculated results. Shorter payback/parity times are calculated when fossil carbon intensity, forest growth rate and biomass conversion efficiency increases, or when the initial forest carbon stock and/or harvest level decreases. Shorter payback/parity times are also calculated when the researcher choose landscape level over stand level carbon accounting (if carbon accounting starts at the harvest rather than at the planting event.) Conversely, longer payback/parity times are calculated when carbon intensity, growth rate and conversion efficiency decreases, or when the initial carbon stock and/or harvest level increases, or the researcher choose stand level over landscape level carbon accounting.
Critics argue that unrealistic system boundary choices are made, or that narrow system boundaries lead to misleading conclusions. Others argue that the wide range of results shows that there is too much leeway available and that the calculations therefore are useless for policy development. EU's Join Research Center agrees that different methodologies produce different results, but also argue that this is to be expected, since different researchers consciously or unconsciously choose different alternative scenarios/methodologies as a result of their ethical ideals regarding man's optimal relationship with nature. The ethical core of the sustainability debate should be made explicit by researchers, rather than hidden away.
Comparisons of GHG emissions at the point of combustion
GHG emissions per produced energy unit at the point of combustion depend on moisture content in the fuel, chemical differences between fuels and conversion efficiencies. For example, raw biomass can have higher moisture content compared to some common coal types. When this is the case, more of the wood's inherent energy must be spent solely on evaporating moisture, compared to the drier coal, which means that the amount of CO2 emitted per unit of produced heat will be higher.
Many biomass-only combustion facilities are relatively small and inefficient, compared to the typically much larger coal plants. Further, raw biomass (for instance wood chips) can have higher moisture content than coal (especially if the coal has been dried). When this is the case, more of the wood's inherent energy must be spent solely on evaporating moisture, compared to the drier coal, which means that the amount of CO2 emitted per unit produced heat will be higher. This moisture problem can be mitigated by modern combustion facilities.
Forest biomass on average produces 10-16% more CO2 than coal. However, focusing on gross emissions misses the point, what counts is the net climate effect from emissions and absorption, taken together. IEA Bioenergy concludes that the additional CO2 from biomass "[...] is irrelevant if the biomass is derived from sustainably managed forests."
Climate impacts expressed as varying with time
The use of boreal stemwood harvested exclusively for bioenergy have a positive climate impact only in the long term, while the use of wood residues have a positive climate impact also in the short to medium term.
Short carbon payback/parity times are produced when the most realistic no-bioenergy scenario is a traditional forestry scenario where "good" wood stems are harvested for lumber production, and residues are burned or left behind in the forest or in landfills. The collection of such residues provides material which "[...] would have released its carbon (via decay or burning) back to the atmosphere anyway (over time spans defined by the biome's decay rate) [...]." In other words, payback and parity times depend on the decay speed. The decay speed depends on a.) location (because decay speed is "[...] roughly proportional to temperature and rainfall [...]"), and b.) the thickness of the residues. Residues decay faster in warm and wet areas, and thin residues decay faster than thick residues. Thin residues in warm and wet temperate forests therefore have the fastest decay, while thick residues in cold and dry boreal forests have the slowest decay. If the residues instead are burned in the no-bioenergy scenario, e.g. outside the factories or at roadside in the forests, emissions are instant. In this case, parity times approach zero.
Like other scientists, the JRC staff note the high variability in carbon accounting results, and attribute this to different methodologies. In the studies examined, the JRC found carbon parity times of 0 to 400 years for stemwood harvested exclusively for bioenergy, depending on different characteristics and assumptions for both the forest/bioenergy system and the alternative fossil system, with the emission intensity of the displaced fossil fuels seen as the most important factor, followed by conversion efficiency and biomass growth rate/rotation time. Other factors relevant for the carbon parity time are the initial carbon stock and the existing harvest level; both higher initial carbon stock and higher harvest level means longer parity times. Liquid biofuels have high parity times because about half of the energy content of the biomass is lost in the processing.
Climate impacts expressed as static numbers
EU's Joint Research Centre has examined a number of bioenergy emission estimates found in literature, and calculated greenhouse gas savings percentages for bioenergy pathways in heat production, transportation fuel production and electricity production, based on those studies. The calculations are based on the attributional LCA accounting principle. It includes all supply chain emissions, from raw material extraction, through energy and material production and manufacturing, to end-of-life treatment and final disposal. It also includes emissions related to the production of the fossil fuels used in the supply chain. It excludes emission/absorption effects that takes place outside its system boundaries, for instance market related, biogeophysical (e.g. albedo), and time-dependent effects. The authors conclude that "[m]ost bio-based commodities release less GHG than fossil products along their supply chain; but the magnitude of GHG emissions vary greatly with logistics, type of feedstocks, land and ecosystem management, resource efficiency, and technology."
Because of the varied climate mitigation potential for different biofuel pathways, governments and organizations set up different certification schemes to ensure that biomass use is sustainable, for instance the RED (Renewable Energy Directive) in the EU and the ISO standard 13065 by the International Organization for Standardization. In the US, the RFS (Renewables Fuel Standard) limit the use of traditional biofuels and defines the minimum life-cycle GHG emissions that are acceptable. Biofuels are considered traditional if they achieve up to 20% GHG emission reduction compared to the petrochemical equivalent, advanced if they save at least 50%, and cellulosic if the save more than 60%.
The EU's Renewable Energy Directive (RED) states that the typical greenhouse gas emissions savings when replacing fossil fuels with wood pellets from forest residues for heat production varies between 69% and 77%, depending on transport distance: When the distance is between 0 and 2500 km, emission savings is 77%. Emission savings drop to 75% when the distance is between 2500 and 10 000 km, and to 69% when the distance is above 10 000 km. When stemwood is used, emission savings varies between 70% and 77%, depending on transport distance. When wood industry residues are used, savings varies between 79% and 87%.
Since the long payback and parity times calculated for some forestry projects is seen as a non-issue for energy crops (except in the cases mentioned above), researchers instead calculate static climate mitigation potentials for these crops, using LCA-based carbon accounting methods. A particular energy crop-based bioenergy project is considered carbon positive, carbon neutral or carbon negative based on the total amount of CO2 equivalent emissions and absorptions accumulated throughout its entire lifetime: If emissions during agriculture, processing, transport and combustion are higher than what is absorbed (and stored) by the plants, both above and below ground, during the project's lifetime, the project is carbon positive. Likewise, if total absorption is higher than total emissions, the project is carbon negative. In other words, carbon negativity is possible when net carbon accumulation more than compensates for net lifecycle greenhouse gas emissions.
Typically, perennial crops sequester more carbon than annual crops because the root buildup is allowed to continue undisturbed over many years. Also, perennial crops avoid the yearly tillage procedures (plowing, digging) associated with growing annual crops. Tilling helps the soil microbe populations to decompose the available carbon, producing CO2.
There is now (2018) consensus in the scientific community that "[...] the GHG [greenhouse gas] balance of perennial bioenergy crop cultivation will often be favourable [...]", also when considering the implicit direct and indirect land use changes.
Albedo and evapotranspiration
Environmental impacts
The environmental impacts of biomass production need to be taken into account. For instance, in 2022, IEA stated that "bioenergy is an important pillar of decarbonisation in the energy transition as a near zero-emission fuel", and that "more efforts are needed to accelerate modern bioenergy deployment to get on track with the Net Zero Scenario [....] while simultaneously ensuring that bioenergy production does not incur negative social and environmental consequences."
Sustainable forestry and forest protection
IPCC states that there is disagreement about whether the global forest is shrinking or not, and quote research indicating that tree cover has increased 7.1% between 1982 and 2016. The IPCC writes: "While above-ground biomass carbon stocks are estimated to be declining in the tropics, they are increasing globally due to increasing stocks in temperate and boreal forests [...]."
Old trees have a very high carbon absorption rate, and felling old trees means that this large potential for future carbon absorption is lost. There is also a loss of soil carbon due to the harvest operations.
Old trees absorb more CO2 than young trees because of the larger leaf area in full-grown trees. However, the old forest (as a whole) will eventually stop absorbing CO2 because CO2 emissions from dead trees cancel out the remaining living trees' CO2 absorption. The old forest (or forest stands) are also vulnerable for natural disturbances that produces CO2. The IPCC found that "[...] landscapes with older forests have accumulated more carbon but their sink strength is diminishing, while landscapes with younger forests contain less carbon but they are removing CO2 from the atmosphere at a much higher rate [...]."
The IPCC states that the net climate effect from conversion of unmanaged to managed forest can be positive or negative, depending on circumstances. The carbon stock is reduced, but since managed forests grow faster than unmanaged forests, more carbon is absorbed. Positive climate effects are produced if the harvested biomass is used efficiently. There is a tradeoff between the benefits of having a maximized forest carbon stock, not absorbing any more carbon, and the benefits of having a portion of that carbon stock "unlocked", and instead working as a renewable fossil fuel replacement tool, for instance in sectors which are difficult or expensive to decarbonize.
The "competition" between locked-away and unlocked forest carbon might be won by the unlocked carbon: "In the long term, using sustainably produced forest biomass as a substitute for carbon-intensive products and fossil fuels provides greater permanent reductions in atmospheric CO2 than preservation does."
IEA Bioenergy writes: "forests managed for producing sawn timber, bioenergy and other wood products can make a greater contribution to climate change mitigation than forests managed for conservation alone." Three reasons are given:
reducing ability to act as a carbon sink when the forest matures.
Wood products can replace other materials that emitted more GHGs during production.
"Carbon in forests is vulnerable to loss through natural events such as insect infestations or wildfires"
Data from FAO show that most wood pellets are produced in regions dominated by sustainably managed forests, such as Europe and North America. Europe (including Russia) produced 54% of the world's wood pellets in 2019, and the forest carbon stock in this area increased from 158.7 to 172.4 Gt between 1990 and 2020. In the EU, above-ground forest biomass increases with 1.3% per year on average, however the increase is slowing down because the forests are maturing.
United Kingdom Emissions Trading System allows operators of CO2 generating installations to apply zero emissions factor for the fraction used for non-energy purposes, while energy purposes (electricity generation, heating) require additional sustainability certification on the biomass used.
Biodiversity
Biomass production for bioenergy can have negative impacts on biodiversity. Oil palm and sugar cane are examples of crops that have been linked to reduced biodiversity. In addition, changes in biodiversity also impacts primary production which naturally effects decomposition and soil heterotrophic organisms.
Win-win scenarios (good for climate, good for biodiversity) include:
Increased use of whole trees from coppice forests, increased use of thin forest residues from boreal forests with slow decay rates, and increased use of all kinds of residues from temperate forests with faster decay rates;
Multi-functional bioenergy landscapes, instead of expansion of monoculture plantations;
Afforestation of former agricultural land with mixed or naturally regenerating forests.
Win-lose scenarios (good for the climate, bad for biodiversity) include afforestation on ancient, biodiversity-rich grassland ecosystems which were never forests, and afforestation of former agricultural land with monoculture plantations.
Lose-win scenarios (bad for the climate, good for biodiversity) include natural forest expansion on former agricultural land.
Lose-lose scenarios include increased use of thick forest residues like stumps from some boreal forests with slow decay rates, and conversion of natural forests into forest plantations.
Pollution
Other problems are pollution of soil and water from fertiliser/pesticide use, and emission of ambient air pollutants, mainly from open field burning of residues.
The traditional use of wood in cook stoves and open fires produces pollutants, which can lead to severe health and environmental consequences. However, a shift to modern bioenergy contribute to improved livelihoods and can reduce land degradation and impacts on ecosystem services. According to the IPCC, there is strong evidence that modern bioenergy have "large positive impacts" on air quality. Traditional bioenergy is inefficient and the phasing out of this energy source has both large health benefits and large economic benefits. When combusted in industrial facilities, most of the pollutants originating from woody biomass reduce by 97-99%, compared to open burning. Combustion of woody biomass produces lower amounts of particulate matter than coal for the same amount of electricity generated.
| Technology | Fuel | null |
2467081 | https://en.wikipedia.org/wiki/Nosean | Nosean | Nosean, also known as noselite, is a mineral of the feldspathoid group with formula: Na8Al6Si6O24(SO4). H2O. It forms isometric crystals of variable color: white, grey, blue, green, to brown. It has a Mohs hardness of 5.5 to 6 and a specific gravity of 2.3 to 2.4. It is fluorescent. It is found in low-silica igneous rocks. There is a solid solution between nosean and hauyne, which contains calcium.
It was first described in 1815 from the Rhineland in Germany and named after the German mineralogist K. W. Nose (1753–1835). The mineral is rare but widespread, found in such diverse localities as ocean islands (e.g., Tahiti) and the La Sal Range in Utah.
| Physical sciences | Silicate minerals | Earth science |
2467302 | https://en.wikipedia.org/wiki/Sodium%20dithionite | Sodium dithionite | Sodium dithionite (also known as sodium hydrosulfite) is a white crystalline powder with a sulfurous odor. Although it is stable in dry air, it decomposes in hot water and in acid solutions.
Structure
The structure has been examined by Raman spectroscopy and X-ray crystallography. The dithionite dianion has C symmetry, with almost eclipsed with a 16° O-S-S-O torsional angle. In the dihydrated form (), the dithionite anion has gauche 56° O-S-S-O torsional angle.
A weak S-S bond is indicated by the S-S distance of 239 pm, which is elongated by ca. 30 pm relative to a typical S-S bond. Because this bond is fragile, the dithionite anion dissociates in solution into the [SO2]− radicals, as has been confirmed by EPR spectroscopy. It is also observed that 35S undergoes rapid exchange between S2O42− and SO2 in neutral or acidic solution, consistent with the weak S-S bond in the anion.
Preparation
Sodium dithionite is produced industrially by reduction of sulfur dioxide. Approximately 300,000 tons were produced in 1990. The route using zinc powder is a two-step process:
2SO2 + Zn → ZnS2O4
ZnS2O4 + 2NaOH → Na2S2O4 + Zn(OH)2
The sodium borohydride method obeys the following stoichiometry:
NaBH4 + 8NaOH + 8SO2 → 4Na2S2O4 + NaBO2 + 6H2O
Each equivalent of H− reduces two equivalents of sulfur dioxide. Formate has also been used as the reductant.
Properties and reactions
Hydrolysis
Sodium dithionite is stable when dry, but aqueous solutions deteriorate due to the following reaction:
2 S2O42− + H2O → S2O32− + 2 HSO3−
This behavior is consistent with the instability of dithionous acid. Thus, solutions of sodium dithionite cannot be stored for a long period of time.
Anhydrous sodium dithionite decomposes to sodium sulfate and sulfur dioxide above 90 °C in the air. In absence of air, it decomposes quickly above 150 °C to sodium sulfite, sodium thiosulfate, sulfur dioxide and trace amount of sulfur.
Redox reactions
Sodium dithionite is a reducing agent. At pH 7, the potential is -0.66 V compared to the normal hydrogen electrode. Redox occurs with formation of bisulfite:
S2O42- + 2 H2O → 2 HSO3− + 2 e− + 2 H+
Sodium dithionite reacts with oxygen:
Na2S2O4 + O2 + H2O → NaHSO4 + NaHSO3
These reactions exhibit complex pH-dependent equilibria involving bisulfite, thiosulfate, and sulfur dioxide.
With organic carbonyls
In the presence of aldehydes, sodium dithionite reacts either to form α-hydroxy-sulfinates at room temperature or to reduce the aldehyde to the corresponding alcohol above a temperature of 85 °C. Some ketones are also reduced under similar conditions.
Uses
Industry
Sodium dithionite is used as a water-soluble reducing agent in some industrial dyeing processes. In the case of sulfur dyes and vat dyes, an otherwise water-insoluble dye can be reduced into its water-soluble alkali metal leuco salt. Indigo dye is sometimes processed in this way.
Domestic and hobby uses
Sodium dithionite can also be used for water treatment, aquarium water conditioners, gas purification, cleaning, and stripping.In addition to the textile industry, this compound is used in industries concerned with leather, foods, polymers, photography, and many others, often as a decolourising agent. It is even used domestically as a decoloring agent for white laundry, when it has been accidentally stained by way of a dyed item slipping into the high temperature washing cycle. It is usually available in 5 gram sachets termed hydrosulfite after the antiquated name of the salt.
It is the active ingredient in "Iron Out Rust Stain Remover", a commercial rust product.
Laboratory
Sodium dithionite is often used in physiology experiments as a means of lowering solutions' redox potential (Eo' -0.66 V vs SHE at pH 7). Potassium ferricyanide is usually used as an oxidizing chemical in such experiments (Eo' ~ .436 V at pH 7). In addition, sodium dithionite is often used in soil chemistry experiments to determine the amount of iron that is not incorporated in primary silicate minerals. Hence, iron extracted by sodium dithionite is also referred to as "free iron."
Aqueous solutions of sodium dithionite were once used to produce 'Fieser's solution' for the removal of oxygen from a gas stream. Pyrithione can be prepared in a two-step synthesis from 2-bromopyridine by oxidation to the N-oxide with a suitable peracid followed by substitution using sodium dithionite to introduce the thiol functional group.
Photography
It is used in Kodak fogging developer, FD-70. This is used in the second step in processing black and white positive images, for making slides. It is part of the Kodak Direct Positive Film Developing Outfit.
Safety
The wide use of sodium dithionite is attributable in part to its low toxicity at 2.5 g/kg (rats, oral).
| Physical sciences | Sulfuric oxyanions | Chemistry |
2467838 | https://en.wikipedia.org/wiki/Blattodea | Blattodea | Blattodea is an order of insects that contains cockroaches and termites. Formerly, termites were considered a separate order, Isoptera, but genetic and molecular evidence suggests they evolved from within the cockroach lineage, cladistically making them cockroaches as well. The Blattodea and the mantis (order Mantodea) are now all considered part of the superorder Dictyoptera. Blattodea includes approximately 4,400 species of cockroach in almost 500 genera, and about 3,000 species of termite in around 300 genera.
Termites are pale-coloured, soft-bodied eusocial insects that live in colonies, whereas cockroaches are darker-coloured (often brown), sclerotized, segmented insects. Within the colony, termites have a caste system, with a pair of mature reproductives, the king and the queen, and numerous sterile workers and soldiers. Cockroaches are not colonial but do have a tendency to aggregate and may be considered pre-social, as all adults are capable of breeding. Other similarities between the two groups include various social behaviours, trail following, kin recognition, and methods of communication.
Phylogeny and evolution
Cladistic analysis of five DNA sequences in 107 species representing all the termite subfamilies, all six cockroach families, including 22 of the 29 subfamilies, and five of the 15 mantis families (as out-groups) showed that the termites are nested within the cockroaches, and that the Cryptocercidae is a sister group to the termites. The mantids were shown to be the sister group to Blattodea. Cryptocercus also shares characteristics such as species of gut bacteria with the termites.
The cockroach families Lamproblattidae and Tryonicidae are not shown but are placed within the superfamily Blattoidea. The cockroach families Corydiidae and Ectobiidae were previously known as the Polyphagidae and Blattellidae.
The evolutionary relationships of the Blattodea (cockroaches and termites), based on Eggleton, Beccaloni & Inward (2007) and modified by Evangelista et al. 2019, are shown in the cladogram:
The cladogram shows the family Alienopteridae (originally assigned to its own order "Alienoptera") as sister to Mantodea; while it was reassigned to the extinct Blattodea superfamily Umenocoleoidea by Vršanský et al., a more recent analysis places Alienopteridae and Umenocoleidae as sister taxa within Dictyoptera, and not within Blattodea.
Diversity
Over 4,000 species of cockroaches are found in every corner of the globe with each continent having its own indigenous species. Most of these are omnivores or detritivores and live in a range of habitats such as among leaf litter, in rotting wood, in thick vegetation, in crevices, in cavities beneath bark, under logs and among debris. Some are arboreal, some live in caves and some are aquatic. A small number of species have taken to living in close proximity to humans in buildings, have been transported around the world by them, and are regarded as pests. Although some species harbour symbionts in their guts which facilitate cellulose digestion, many species also produce enzymes to digest cellulose independent of the symbionts.
Over 3,000 species of termite are found in all the continents except Antarctica. The greatest diversity is found in Africa and relatively few species inhabit Europe and North America. They are also detrivores and many species eat wood, having specialised guts with symbiotic protozoa to digest the cellulose. Termites have soft bodies and keep out of sight as far as possible. They can loosely be subdivided into dampwood, drywood and subterranean types. In general, dampwood termites inhabit coniferous forests, drywood termites inhabit hardwood forests and subterranean termites live in a wide variety of habitats.
Characteristics
Termites are eusocial insects that live in colonies. They have a caste system, with a king and queen in each colony and many non-reproductive workers. The workers forage for food which they bring back to the colony to feed the reproductives and the developing young. Cockroaches are also social insects but do not live in colonies, and all adults are able to reproduce. Some species form aggregations, others show an inclination to aggregate, and some exhibit parental care of their offspring.
Cockroaches and termites have striking similarities in behaviour which they likely inherited from their common ancestor. These include an attraction to warm and humid places, thigmotaxis, burrowing, substrate manipulation, hygienic behaviour, food sharing, cannibalism, excretion behaviour, vibrational communication, kin recognition, trail following, allogrooming, care of the brood, cropping of antennae and certain mating behaviours. In some of these behaviours, there are marked similarities between termites and juvenile, but not adult, cockroaches. During the evolution of eusociality, the individuals need to share a desire to group together. Juvenile cockroaches have a tendency to aggregate while adults often compete aggressively with each other for space and resources. Similarly, grooming and being groomed is common in termite colonies but allogrooming is not a behaviour generally engaged in by cockroaches although individuals groom themselves. An exception to this is the cockroach Cryptocercus, which seems to be more closely related to the termites than to other cockroaches.
Here juveniles groom each other and also groom adults.
Both groups are also affected by their social environments. A single termite, kept alone, has a significantly decreased level of vigour and a shorter lifespan than when two are kept together. An isolated cockroach nymph may grow at less than half the rate of grouped individuals, and has a poorer life expectancy.
Both termites and cockroaches engage in coprophagy, the consumption of fecal pellets. Adult termite workers forage and bring food back to the nest where they pass it to the reproductives and young either by mouth or by anus, providing the whole of their nutritional needs in this manner. Young cockroaches are ineffective foragers, seldom straying from their hiding places, and obtain much of their nourishment from eating the fecal pellets of larger individuals. From these they acquire the microbial flora that helps them to digest their food.
A single cockroach family, the Cryptocercidae, and one primitive species of termite, Mastotermes darwiniensis, share such characteristics as the segmental origin of certain female reproductive structures, and the fact that both deposit their eggs in the oothecae that are typical of cockroaches.
Cockroaches
Arthropods similar to living cockroaches dominated the insect communities of the Carboniferous period. Modern crown group cockroaches radiated from them by the middle of the Mesozoic, with the first appearance of the extant family Corydiidae during the Late Jurassic. This group of insects are nocturnal, only foraging for food and water at night. They are not considered eusocial because their populations are not divided into different caste systems; however, they are still social creatures and can live in groups with over a million individuals. The cockroach is flattened dorsolaterally and is roughly oval with a shield-like plate, the pronotum, covering its thorax and posterior region of the head. The antennae are many-segmented, long and slender, and the mouthparts are adapted for chewing. The forewings are normally leathery and the hind wings membranous. The coxae of the legs are flattened to enable the femurs to fit neatly against them when folded. Cockroaches are hemimetabolous; there is no pupal stage and the nymphs resemble the adults apart from their size and the absence of wings. Female cockroaches produce an egg sac known as an ootheca and can hold anywhere from 12-25 eggs depending upon the species. Some species display parenting behavior, whereas other species have nothing to do with the young. In most species, growth to maturity takes three to four months, but in a few species, the nymph stage can last for several years. The main factors affecting the duration of the nymph stage are seasonal differences, and the amount of nutrients received in the diet.
Chemical communication
As in most insect species, cockroaches communicate with one another by the release of pheromones. It has also been discovered that cockroaches release hydrocarbons from their body that are transferred through interactions of the antennae. These hydrocarbons can aid in cockroach communication and can even tell whether an individual is a member of its kin or not to prevent inbreeding. Cockroaches that have been isolated in a lab setting have shown extreme behavioral effects and are less stimulated by these hydrocarbons and pheromones, possibly suggesting a group environment is required for development of these communication skills.
Termites
All species of termite are to some degree eusocial, and the members of a colony are differentiated into caste systems. The majority of termite populations consist of the worker caste, which are responsible for foraging, nest building, grooming, and brood care. The soldier caste has one responsibility, which is to protect the nest from predators and other competitors. Soldiers have highly developed mandibles as well as many exocrine glands that can secrete multiple defensive substances harmful to predators.
Normally, only the king and queen termite reproduce; the other castes are all sterile. There are two classes of reproductives: primary reproductives and neotenic reproductives. The primary reproductives class is responsible for colony creation and is characterized by compound eyes, wing marks (spots where wings once were before shedding), and defined sclerotization. Neotenic reproductives can develop from within the colony usually when one of the primary reproductives has died, or can develop in addition to the queen. neotenic reproductives can experience two different phenotypes, one with wings and one without. If neotenics are winged they will fly away from the parental colony, pair up and form a new colony, but if they are wingless they will remain within the parental colony. The different developmental routes taken by these two morphs are usually dependent upon food availability in the colony, or varying levels of parasitism within the colony. The caste into which any particular nymph will develop begins to become apparent among the late instars; at this time, potential reproductives will begin to show an increase in the size of the gonadal region.
Termite colonies may be arboreal, mound-like or subterranean, with primitive termites nesting completely inside enclosed structures such as stumps or logs. Nest construction is largely from the termites' own faecal matter, other materials being chewed vegetable fibre, which makes a weak carton-like but waterproof substance, and soil, which makes a strong substance, but which is subject to erosion by water. Aerial nests are connected to the ground by enclosed passageways; the soft-bodied, blind workers of most species live permanently in their protected environments and do not venture into the open air. Trinervitermes trinervoides is an exception to this, with workers foraging in small groups on the surface at night, secreting noxious terpenes to deter predators. The nests are complex structures, and tunnels link them to the foraging areas. In Africa, termite mounds can be as large as nine meters tall and thirty meters in diameter, producing an area of increased fertility and creating a small hotspot for biodiversity.
| Biology and health sciences | Cockroaches & Termites (Blattodea) | Animals |
2467853 | https://en.wikipedia.org/wiki/Nothosaurus | Nothosaurus | Nothosaurus ('false lizard', from the Ancient Greek and ) is an extinct genus of sauropterygian reptile from the Triassic period, approximately 240–210 million years ago, with fossils being distributed throughout the former Tethys Ocean, from North Africa and Europe to China. It is the best known member of the nothosaur order.
Description
Nothosaurus was a semi-oceanic animal which probably had a lifestyle similar to that of today's seals. It was about , with long, webbed toes and possibly a fin on its tail. However, some species such as N. zhangi and N. giganteus were larger, up to . When swimming, Nothosaurus would use its tail, legs, and webbed feet to propel and steer it through the water. The skull was broad and flat, with long jaws, lined with needle teeth, it probably caught fish and other marine creatures. Trackways attributed, partly by process of elimination, to a nothosaur, that were reported from Yunnan, China in June 2014, were interpreted as the paddle impressions left as the animals dug into soft seabed with rowing motions of their paddles, churning up hidden benthic creatures that they snapped up. Once caught, few animals would be able to shake themselves free from the mouth of Nothosaurus.
In many respects its body structure resembled that of the much later plesiosaurs, but it was not as well adapted to an aquatic environment. It is thought that one branch of the nothosaurs may have evolved into pliosaurs such as Liopleurodon, a short-necked plesiosaur that grew up to , and the long-necked Cryptoclidus, a fish eater with a neck as long as .
Species
There are nearly a dozen known species of Nothosaurus. The type species is N. mirabilis, named in 1834 from the Germanic Muschelkalk. Other species include N. giganteus (previously known as Paranothosaurus) from Osnabrück, Germany; N. juvenilis, also from Germany; N. edingerae from the Upper Muschelkalk and Lower Keuper; N. haasi and N. tchernovi from Makhtesh Ramon, Israel; N. cymatosauroides from the Spanish Muschelkalk; N. jagisteus from the Upper Muschelkalk of Hohenlohe, Germany; and N. youngi, N. yangjuanensis (and its junior synonym N. rostellatus) and the recently named N. zhangi from Guizhou, China. Several species have been described from the Lower Muschelkalk in Winterswijk, the Netherlands, including N. marchicus (and its junior synonym N. winterswijkensis) and N. winkelhorsti. Recently, the long considered lost type material of N. schimperi Meyer, 1842 from the Lower Muschelkalk of Soultz-les-Bains, Alsace, France, has been rediscovered and a lectotype has been designated.
Klein and Albers (2009) conducted a phylogenetic analysis, but did not test the monophyly of Nothosaurus, as other nothosaurids were not included in their analysis.
Several other species have been named but are now generally considered invalid. One such species, N. procerus, is now considered a junior subjective synonym of N. marchicus. Other species now considered junior synonyms of N. marchicus include N. crassus, N. oldenburgi, N. raabi, N. schroderi, N. venustus and the recently named N. winterswijkensis. Junior synonyms of N. giganteus, the second largest Nothosaurus species, include N. andriani, N. angustifronis, N. aduncidens, N. baruthicus and N. chelydrops.
A species level phylogenetic analysis of Nothosauridae was performed by Liu et al. (2014), and included all known valid species of the family and Nothosaurus apart from Lariosaurus stensioi (type of Micronothosaurus), Nothosaurus cymatosauroides, and Ceresiosaurus lanzi. Due to the inclusion of other nothosaurids other than Nothosaurus, the monophyly of Nothosaurus was tested for the first time. The analysis found both Lariosaurus and Nothosaurus to be polyphyletic in regard to each other and all the other genera of the family, making a systematic revision of these two genera necessary. Below, their results are shown with type species of named nothosaurid genera noted. Later, in 2017, the species N. juvenilis, N. youngi, and N. winkelhorsti were formally moved to Lariosaurus.
| Biology and health sciences | Prehistoric marine reptiles | Animals |
2469649 | https://en.wikipedia.org/wiki/Tylosaurus | Tylosaurus | {{Automatic taxobox
| fossil_range = Turonian-Maastrichtian,~
| image = Bunker Tylosaur.png
| image_upright = 1.15
| image_caption = Mounted cast of the T. proriger "Bunker" specimen (KUVP 5033)
| taxon = Tylosaurus
| authority = Marsh, 1872
| type_species = Tylosaurus proriger
| type_species_authority = (Cope, 1869)
| subdivision_ranks = Other species
| subdivision = * T. nepaeolicus (Cope, 1874)
T. bernardi (Dollo, 1885)
T. gaudryi (Thevenin, 1896)
T. ivoensis (Persson, 1963)
T. iembeensis (Antunes, 1964)
T. pembinensis (Nicholls, 1988)
T. saskatchewanensis Jiménez-Huidobro et al., 2018
{{collapsible list|title=Disputed or unpublished|
T. kansasensis Everhart, 2005
T. "borealis" Garvey, 2020
}}
| synonyms =
}}Tylosaurus (; "knob lizard") is a genus of russellosaurine mosasaur (an extinct group of predatory marine lizards) that lived about 92 to 66 million years ago during the Turonian to Maastrichtian stages of the Late Cretaceous. Its fossils have been found primarily around North Atlantic Ocean including in North America, Europe, and Africa.
Research history
Possible first finds
The earliest Tylosaurus fossils were likely discovered by various Native American peoples and may have been the source of much of their folklore, with the earliest known ones dating back to well before the arrival of European settlers, around the 1500s. More recent accounts from peoples living in the Great Plains even speak of an ancient era ruled by massive aquatic creatures that were in constant combat with thunderbirds and were petrified by them. The considerable presence of fossils of large mosasaurs such as Tylosaurus and pterosaurs such as Pteranodon in this region may have been the origins of these myths.
In 1804, the Lewis and Clark Expedition discovered a now-lost fossil skeleton alongside the Missouri River, which was identified as a long fish. In 2003, Richard Ellis speculated that the remains may have belonged to Mosasaurus missouriensis. Alternatively, a 2007 study led by Robert W. Meredith and colleagues suggested that the fossils would possibly come from a tylosaurine mosasaur based on the measurements cited by Clark and Gass and the evidence of Tylosaurus fossils that have been found in the Missouri River. However, the authors also mentioned the possibility that the remains would also come from an elasmosaurid plesiosaur, which are also known from the river, although being rarer.
First formal discoveries
Tylosaurus was the third new genus of mosasaur to be described from North America behind Clidastes and Platecarpus and the first in Kansas. The early history of the genus as a taxon was subject to complications spurred by the infamous rivalry between American paleontologists Edward Drinker Cope and Othniel Charles Marsh during the Bone Wars. The type specimen was described by Cope in 1869 based on a fragmentary skull measuring nearly in length and thirteen vertebrae lent to him by Louis Agassiz of the Harvard Museum of Comparative Zoology. The fossil, which remains in the same museum under the catalog number MCZ 4374, was recovered from a deposit of the Niobrara Formation located in the vicinity of Monument Rocks near the Union Pacific Railroad at Fort Hays, Kansas. Cope's first publication of the fossil was very brief and was named Macrosaurus proriger, the genus being a preexisting European mosasaur taxon. The specific epithet proriger means "prow-bearing", which is in reference to the specimen's unique prow-like elongated rostrum and is derived from the Latin word prōra (prow) and suffix -gero (I bear). In 1870, Cope published a more thorough description of MCZ 4374. Without explanation, he moved the species into another European genus Liodon and declared his original Macrosaurus proriger a synonym.
In 1872, Marsh argued that Liodon proriger is taxonomically distinct from the European genus and must be assigned a new one. For this, he erected the genus Rhinosaurus, which means "nose lizard" and is a portmanteau derived from the Ancient Greek words (, meaning "nose") and (, meaning "lizard"). Cope responded by arguing that Rhinosaurus was already a preoccupied synonym of Liodon. He disagreed with Marsh's arguments but proposed that in case Marsh was indeed correct, the genus name Rhamphosaurus should be used. Marsh later discovered that the taxon Rhamphosaurus was preoccupied as a genus of lizard named in 1843. As a result, he suggested a move to a newly erected genus named Tylosaurus. This name means "knob lizard" in another reference to the elongated rostrum characteristic of the genus. It is derived from the Latin tylos (knob) and Ancient Greek . Despite coining the new genus, Marsh never formally transferred this Rhinosaurus species to Tylosaurus; this was first done in 1873 by Joseph Leidy. Tylosaurus subsequently became the almost universally accepted genus to include this species, the exception to this adoption being Cope, who refused to accept Marsh's new genus and continued to refer to its species as Liodon. Cope's persistence can be seen in his 1874 description of another species of Tylosaurus, which he named Liodon nepaeolicus. The type specimen of this species was discovered by geologist Benjamin Franklin Mudge near the Solomon River, and consists of several cranial fragments and a dorsal vertebra now catalogued as AMNH 1565. This species, whose specific epithet refers to Nepaholla, the Native American name for the Salomon River is formally transferred to the genus Tylosaurus in 1894 by John Campbell Merriam.
Later discoveries and other species
In his major work published in 1967, Dale A. Russell recognized only two valid species in Tylosaurus, namely T. proriger and T. nepaeolicus. However, throughout the 19th and 20th centuries, many species of mosasaurs coming from around the world, originally described as being from separate genera, were now recognized as belonging to Tylosaurus.
In 1885, Louis Dollo described the genus and species Hainosaurus bernardi from an almost complete but poorly preserved skeleton discovered in a phosphate quarry in the Ciply Basin in Belgium, the specimen having since been catalogued as IRSNB R23. The prefix Haino- in the generic name refers to the Haine, a river located nearby the Ciply Basin, and thus combined with means "lizard from the Haine". The specific epithet is named in honor of Leopold Bernard, who made the excavation and exhumation of the specimen possible. In 1988, a second species historically pertained to Hainosaurus was described by Elizabeth Nicholls based on a partial skeleton catalogued as MT 2 and having been discovered in Manitoba, Canada. The specific epithet refers to the type locality of the taxon, namely the Pembina Member of the Pierre Shale. The attribution of H. pembinensis to Hainosaurus is first discussed by Johan Lindgren in 2005, but it was in a revision published in 2010 that the species was moved to Tylosaurus by Timon Bullard and Michael Caldwell, being then renamed as T. pembinensis. In this same revision, the authors suggested that a redescription of the type species H. bernardi would be necessary in order to know if Hainosaurus should be maintained as a distinct genus. This redescription was finally carried out by Paulina Jimenez-Huidobro and Caldwell in 2016, in which they transferred the species to Tylosaurus, being then renamed as T. bernardi. Although this new combination has been widely recognized since, some authors nevertheless suggest continuing to maintain the genus Hainosaurus as distinct, justified in particular on the basis of dental traits not detailed in the 2016 revision.
The fourth recognized species of the genus was described in 1896 by Armand Thevenin on the basis of a partial skull discovered at Éclusier-Vaux, in Somme, France. In his description, Thevenin thinks that this specimen, since catalogued as MNHN 1896–15, would represent a species of Mosasaurus, initially naming it Mosasaurus gaudryi. The specific epithet is named in honor of his mentor Jean Albert Gaudry, the latter having previously studied the skull and thinking that it would come from a species of Liodon. In 1992, Theagarten Lingham-Soliar reassigned the species to Hainosaurus, the latter having previously interpreted the holotype of this taxon as an additional specimen of H. bernardi. In 2005, Lindgren moved this species to Tylosaurus, notably due to its dental characteristics being closer to other lineages of the genus.
In 1963, Per Ove Persson identified a new mosasaurid on the basis of isolated teeth discovered in a deposit located in an area called Ivö Klack, near Ivö Lake in the Kristianstad Basin in Scania, Sweden. Fossils from this same mosasaurid have been documented in this area since 1836, but it is from that year onwards that they are described as coming from a subspecies of Mosasaurus hoffmannii, being then named M. hoffmannii ivoensis, the second specific epithet referring to the type locality. In 1967, Russell elevated the taxon to a separate species within the genus, and assigned to it fossils from the Niobrara Formation of Kansas, including a partial skull. When the taxon was significantly revised in a in 2002 study, being reassigned to Tylosaurus, Lindgren and Mikael Siverson referred additional fossils to this latter that had been discovered at Ivö Klack, including cranial and vertebral remains. In their study, the authors also found that Russell's attributions of the Kansas fossils to this species were erroneous, the remains coming from a distinct taxon. In a 2008 paper, Caldwell and colleagues suggested that T. ivoensis might belong in the related genus Taniwhasaurus based on its dental features, but the authors see this as a subject for another study.
In 1964, Miguel Telles Antunes described the species Mosasaurus iembeensis from a partial skull excavated from the Itombe Formation near the town of Iembe (hence the name), Angola. In 1992, Lingham-Soliar argued that the cranial features were not consistent with those of Mosasaurus and were more characteristic of Tylosaurus, the species being renamed as T. iembeensis. However, the author did not identify the holotype skull, which he considered to reside in the collections of the NOVA University Lisbon without a catalogue number, and it is since 2006 reported as being destroyed in a fire. In 2012, Octávio Mateus and colleagues reported that an additional specimen of T. iembeensis consisting of fragmentary cranial elements was recovered during an expedition to the locality of the since-lost holotype, although it was not figured or formally described.
In 2005, Michael J. Everhart described the species T. kansasensis based on several specimens that had been discovered in Kansas, again in the fossil record of the Niobrara Formation. The holotype specimen consists of a well-preserved skull and six cervical vertebrae cataloged as FHSM VP-2295, which was discovered in 1968 in Ellis County. The validity of this species was questioned as early as 2007 by Caldwell, to the point that in a study published in 2016 with Jiménez-Huidobro and other authors, the latter considers it a juvenile form of T. nepaeolicus, thus making the first name a junior synonym of the second. This is disputed by Everhart himself in a 2017 book, but he only comments on the study as "poorly researched and written" without detailing how. An ontogenetic review of Tylosaurus conducted by Robert F. Stewart and Jordan Mallon in 2018 favors maintaining T. kansasensis as valid, while another conducted in 2020 by Amelia R. Zietlow prefers to follow the advice set out in the 2016 revision.
In 2006, Bullard wrote a Master of Science thesis describing the species T. saskatchewanensis from a partial skeleton catalogued as RSM P2588.1. This specimen, nicknamed "Omācīw" (meaning "hunter" in Cree), was discovered in 1994 near Herbert Ferry, at the Lake Diefenbaker, Saskatchewan. Although originally described informally and via incompletely prepared fossils, the proposed taxon was nevertheless recognized as valid in some subsequent studies. In 2018, Bullard co-authored a multi-author study led by Jiménez-Huidobro which formally described Omācīw, which by then was more fully prepared, and confirmed its identity as belonging to a distinct species.
In 2020, Samuel Garvey wrote a thesis on a partial skull of Tylosaurus catalogued as TMP 2014.011.0001. With visibly distinct features from other species and having been discovered approximately northeast of Grande Prairie, Alberta, this makes the specimen the northernmost known occurrence of the genus, being then named T. borealis, in reference to its northernly presence.
Depiction history
When Cope described the holotype specimen of T. proriger in 1870, he visualized it as an "excessively elongated reptile", due to the morphology of the caudal vertebrae which suggested this. Taking into account his descriptions, this would result in a sea serpent-like reptile reaching lengths rivaling those of the largest cetaceans. The following year, Cope added more details to his visualization of the animal. For him, the head of Tylosaurus would be conical in shape, with eyes on top, and having a jaw connected to a throat similar to that of a pelican, thus facilitating the entry of its prey. Still according to Cope, the animal would have had only the flippers located at the front of the body, those at the back being absent. The tail is seen as long and flat, used in eel-like locomotion. This depiction was followed in various works published during the late 19th century, although some depictions also depict the animal with a long neck.
In a major revision of mosasaurs published in 1898, Samuel Wendell Williston provided a new anatomical description of Tylosaurus that corrected many of the misconceptions of earlier paleontologists about the genus. Specifically, his paper included a rigorous skeletal reconstruction of T. proriger based on three partial specimens from the collections of the University of Kansas Natural History Museum. Thus, this reconstruction depicts the animal as very mobile marine predator with four flippers, a short neck and a much shorter tail than previous depictions, Williston also fixing a maximum body measurement close to those still cited today, i. e. long. Despite the fact that the spinal column is drawn as straight and not as curved, this reconstruction is still recognized as valid by the scientific community. The discovery of the first known substantially complete skeleton of Tylosaurus was revealed as early as 1899 by Henry Fairfield Osborn, followed by other more or less similar finds which were made from the beginning of the 1900s.
Description
Tylosaurus was a type of derived mosasaur, or a latecoming member with advanced evolutionary traits such as a fully aquatic lifestyle. As such, it had a streamlined body, an elongated tail ending with a downturn supporting a two-lobed fin, and two pairs of flippers. While in the past derived mosasaurs were depicted as akin to giant flippered sea snakes, it is now understood that they were more similar in build to other large marine vertebrates such as ichthyosaurs, marine crocodylomorphs, and archaeocete whales through convergent evolution.
Size
Some species of Tylosaurus are among the largest known mosasaurs. The largest well-known specimen, a skeleton of T. proriger from the University of Kansas Natural History Museum nicknamed "Bunker" (KUVP 5033), has been estimated to measure between long. A fragmentary skeleton of another T. proriger from the Sternberg Museum of Natural History (FHSM VP-2496) may be from an even larger individual; Everhart estimated the specimen to come from a individual compared to his estimate for Bunker. The genus exhibits Cope's rule, in which its body size has been observed to generally increase over geologic time. In North America, the earliest representatives of Tylosaurus during the Turonian and Coniacian (90-86 mya), which included early T. nepaeolicus and its precursors, typically measured long and weighed between . During the Santonian (86-83 mya), T. nepaeolicus and newly-appearing T. proriger were long and weighed around . By the Early Campanian, T. proriger attained lengths of . Everhart speculated that because mosasaurs continuously grew throughout their lifetime, it would have been possible for some extremely old Tylosaurus individuals to reach in absolute maximum length. However, he stressed the lack of fossil evidence suggesting such sizes and the odds against any being preserved.
Other Campanian-Maastrichtian species were similarly large. The most recent maximum estimate for T. bernardi is by Lindgren (2005); historically the species was erroneously estimated at even larger sizes of . A reconstruction of T. saskatchewanensis by the Royal Saskatchewan Museum estimated a total length of over . A mounted skeleton of T. pembinensis, nicknamed "Bruce," at the Canadian Fossil Discovery Centre measures at long and was awarded a Guinness World Records for "Largest mosasaur on display" in 2014. However, the skeleton was assembled for display prior to Bullard and Caldwell (2010)'s reassessment that found the species' number of vertebrae to be exaggerated. T. "borealis" is estimated at in total length.
Skull
The largest known skull of Tylosaurus is T. proriger KUVP 5033 (the "Bunker" specimen), estimated at long. Depending on age and individual variation, Tylosaurus skulls were between 13 and 14% of the total skeleton length. The head was strongly conical and the snout proportionally longer than most mosasaurs, with the exception of Ectenosaurus.
Cranium
The most recognizable characteristic of Tylosaurus is the elongated edentulous rostrum that protrudes from its snout, for which the genus is named. This is formed by the elongation of the front end of the premaxilla and dentary. The rostrum was small and acutely angled at birth, but rapidly developed into a blunt, elongated "knob." The snout is heavily built, supported by a broad and robust internarial bar (comprising the posterodorsal process of the premaxilla, nasals, and anterior process of the frontal), which provided effective shock absorption and stress transfer. Because of this, it has been proposed that the tylosaurine rostrum was elongated for use in ramming prey or rivals, but recent research on Taniwhasaurus found a complex neurovascular system in the snout, suggesting that the rostrum was extremely sensitive, and therefore it is unlikely that the rostrum was used as a ramming weapon. The snout holds the terminal branches for the trigeminal nerves through randomly scattered foramina on the rostrum and along the ventral margin of the maxilla, above the gum line.
The premaxilla, maxilla, and frontal bones border the external nares, or body nostril openings; unlike other mosasaurs, the prefrontal bones are excluded from the border of the nares by a long posterodorsal process of the maxilla. The nares open above the fourth maxillary tooth anteriorly in T. proriger and T. pembinensis, between the third and fourth tooth in T. nepaeolicus, and posterior to the fourth tooth in T. bernardi. Nare length relative to skull length varied between species: it is proportionally short in T. proriger (20-27% skull length), T. bernardi (24% skull length), and T. gaudryi (25-27% skull length), and long in T. pembinensis (28-31% skull length). The nasal bones were either free-floating or lightly articulated to the internarial bar, did not contact the frontal, and were not fused to each other as they are in extant varanid lizards. The nasals' loose association with the rest of the skull in Tylosaurus and other mosasaurs may be why the bones are frequently lost and therefore exceedingly rare; Tylosaurus is one of the only mosasaurs in which the nasal bones are clearly documented; the other is the holotype of Plotosaurus, although one of the bones is missing.
The external nares lead to the choanae (internal nares) in the palate, which provide passage from the nostrils to the throat. In Tylosaurus, they are shaped like a compressed teardrop and bordered by the vomers, palatines, and the maxilla. Anterior to the choanae, each vomer borders the fenestra for the Jacobson's organ, which is involved in the tongue-based sense of smell. It begins opposite of the fourth maxillary tooth in Tylosaurus, and also ends immediately past the fifth maxillary tooth in T. bernardi. The exit point for the veins leading to sinuses inside the palatine occur right in front of the Jacobson's organ between the vomers and maxilla. This differs from living varanids, where the exit occurs behind the organ.
The frontal bone in Tylosaurus usually, but not always, possesses a low midline crest. It is most prominent in T. proriger, and is moderately developed in T. saskatchewanensis and T. bernardi, extending onto the premaxilla in the latter. The frontal crest is present but poorly developed in most T. nepaeolicus skulls, and occasionally lost in some mature individuals. The frontal overlaps the prefrontals and postorbitofrontals above the orbits (eye sockets), and the parietal posteriorly. The position of the pineal eye on the parietal is variable, either appearing close to the frontoparietal suture or contacting it. The orbits are bordered by the prefrontal, lacrimal, postorbitofrontal, and jugal bones. A diagnostic feature of Tylosaurus is that the prefrontals and postorbitofrontals overlap above the orbits, preventing contribution of the frontal. The jugal forms the bottom of the orbit; in Tylosaurus, it is L-shaped and has a distinctive serif-like extension at the lower back corner of the junction between the horizontal and vertical rami (arms) called the posteroventral process. The vertical ramus is overlapped by the postorbitofrontal in most species, and the horizontal ramus overlaps the maxilla. In T. bernardi, the vertical ramus is not overlapped but joins with the postorbitofrontal by a suture, and is much thicker than the horizontal ramus.
The quadrate bones (homologous to the incus in mammals) are located at the back of the skull, articulating the lower jaw to the cranium and holding the eardrums. The complex anatomy of the bone renders it extremely diagnostic, even to the species level. In lateral view, the quadrate resembles a hook in immature T. nepaeolicus and T. proriger individuals, but in adult forms for both species and in T. bernardi, T. pembinensis, and T. saskatchweanensis it takes on a robust oval-like shape. The eardrum (tympanum) attached to the lateral sufrace of the bone within a bowl-like depression called the alar conch. The conch is shallow in T. nepaeolicus, T. proriger, and T. bernardi, and deep in T. pembinensis and T. saskatchewanensis. The alar rim is thin in T. nepaeolicus, T. proriger, and T. bernardi, and thick in T. bernardi, T. pembinensis, and T. saskatchewanensis. The suprastapedial process is a hook-like extension of bone that curves posteroventrally from the apex of the shaft into an incomplete loop, and it likely served as the attachment point for the depressor mandibulae muscles that opened the lower jaw. The process is slender and proportionally long in immature T. nepaeolicus and T. proriger, and thickened as the animals matured. The process is of similar length to T. proriger in T. saskatchwanensis and shorter in T. bernardi. In T. pembinensis, it abruptly turns medially at a 45° downward angle. A similar deflection appears in some juvenile T. nepaeolicus quadrates. Emerging from the posteroventral margin of the alar conch is the infrastapedial process. Its shape appears to changes ontogenetically in T. nepaeolicus and T. proriger; in the former, the process is absent in juveniles but appears as a small bump in adults, while in T. proriger, it is present as a subtle point in juveniles of and becomes a distinctively broad semicircle in adults. The process is small in T. bernardi, and in T. pembinensis and T. saskatchewanensis, it is rounded. In T. saskatchewanensis, the suprastapedial process almost touches the infrastapedial process. At the bottom of the shaft is the mandibular condyle, which forms the joint between the quadrate and the lower jaw. It is rounded in shape in adults. On the medial surface of the bone, a thick, pillar-like vertical ridge often protrudes beyond the dorsal margin of the quadrate so that it is visible in lateral view.
Jaws and teeth
The upper jaws include the premaxilla and maxilla, and the lower jaws include the dentary, splenial, coronoid, angular, surangular, and prearticluar-articular (like other squamates, the prearticular is fused to the articular). The premaxilla, maxilla, and dentary house the marginal dentition, and the pterygoids house palatal dentition. On each side of the skull, Tylosaurus had 2 premaxillary teeth, 12 to 13 maxillary teeth, 13 dentary teeth, and 10 to 11 pterygoid teeth. The dentition is homodont, meaning that all teeth are nearly identical in size and shape, with the exception of the pterygoid teeth, which are smaller and more recurved than the marginal teeth.
Tylosaurine dentaries were elongate; the dentary is between 56 and 60% of total length of the entire lower jaw in adult T. nepaeolicus and T. proriger, about 55% in T. pembinensis, and 62% in T. saskatchwanensis. The dentary is robust, though not as strongly built as it is in Mosasaurus, Prognathodon, or Plesiotylosaurus. The ventral margin of the dentary ranges from straight to slightly concave. A small dorsal ridge appears anterior to the first dentary tooth in mature individuals of T. proriger.
The marginal dentition of most species is adapted for cutting large marine vertebrates, while those in T. ivoensis and T. gaudryi appear more optimized for piercing or smashing prey, and T. "borealis" in both piercing and cutting. Marginal teeth are triangular with a slight recurve towards the back of the jaws so that the lingual (tongue-facing) side forms a U-shaped curve. From top view, they are compressed at the lingual and labial (lip-facing) sides to form an oval-like shape. Teeth of immature T. proriger are initially compressed, but become conical in adulthood. Carinae (cutting edges) are finely serrated with small denticles except in juvenile T. nepaeolicus. In T. pembinensis, they are faint. The teeth generally have both anterior and posterior carinae, but some anterior teeth may have only anterior carinae. The placement of carinae, if paired, is not always equal; in at least T. proriger, T. ivoensis, T. gaudryi, and T. pembinensis, they are positioned such that the surface area of the tooth's lingual side is greater than the labial side. Both sides are always balanced in area in T. bernardi. The enamel surface is lined with thin fine ridges called striations that run vertically from the tooth's base. The surface is also either smooth or faintly faceted, in which it is flattened into multiple sides to form a prism-like geometry.
Bardet et al. (2006) classified Tylosaurus species into two morphological groups based on marginal dentition. The North American proriger group includes T. proriger and T. nepaeolicus and is characterized by teeth with smooth or faint facets, less prominent carinae, and a vein-like network of primitive striations extending to near the tip. The group was originally defined as having slender teeth, but subsequent research has since recognized that slenderness is an ontogenetic trait in T. proriger with robust teeth appearing in adult forms. Though not formally classified within a group, the marginal teeth of T. saskatchwanensis shares a comparable morphology with T. proriger. The second is the Euro-American ivoensis group and consists of T. ivoensis, T. gaudryi, and T. pembinensis. Their teeth are robust with prominent carinae with striations on the lingual and occasionally labial sides that do not reach the tooth's tip, and facets on the labial side. The facets are gentle in T. pembinensis, while in T. ivoensis they are slightly concave. The latter feature is also known as fluting. Marginal teeth in T. gaudryi are virtually indistinguishable from those in T. ivoensis. T. iembeensis was not placed within either group; no further description is known of its teeth other than having striations and no facets. The distinction of an ivoensis group is contentious. Caldwell et al. (2008) argued that T. pembinensis cannot be compared with T. ivoensis as the former's teeth are not fluted, and that T. ivoensis is more allied with the distinctively fluted teeth of Taniwhasaurus. Jiménez-Huidobro and Caldwell (2019) listed the absence of marginal fluting as a diagnostic (taxon-identifying) trait that differentiates Tylosaurus from Taniwhasaurus.
The pterygoid teeth may have enabled ratchet feeding, in which the upper teeth held prey in place as the lower jaw slides back and forth via a streoptostylic jaw joint. The bases of the pterygoid teeth are nearly circular, and each tooth is divided into front and back-facing sides of near-equal surface area via a pair of faint buccal and lingual carinae, except in T. gaudryi, in which the teeth are mediolaterally compressed. Carinae are not serrated. The anterior surface tends to be either smooth of faintly faceted, while the posterior surface is striated.
Postcranial skeleton
Both pectoral and pelvic girdles are unfused in adult Tylosaurus, in contrast to other taxa (e.g., Prognathodon overtoni). Tylosaurus is also distinguished from other mosasaurs by a scapula that is significantly smaller than the coracoid and the absence of the anterior emargination of the coracoid, as well as the absence of a well-developed pubic tubercle.
Tylosaurus limbs are primitive relative to other mosasaurs; their stylopodia (humeri and femora) lack both the complex muscle attachment sites and extreme proximodistal shortening present in other derived taxa. Both carpals and tarsals in tylosaurines are mostly unossified; while other mosasaurs typically have between three and five carpals and tarsals, adult Tylosaurus never possess more than two ossified carpal bones (usually only the ulnare, sometimes the ulnare and distal carpal four) and two ossified tarsal bones (usually only the astragalus, sometimes the astragalus and distal tarsal four). Hyperphalangy (increased number of phalanges relative to the ancestral condition) is present in both fore- and hindlimbs, and the phalanges are spindle-shaped, unlike the short, blocky hourglass-shaped phalanges possessed by mosasaurines. The pisiform appears to be either unossified or absent in tylosaurines. The functional consequences of differences in limb anatomy across different mosasaur clades is unclear.
Tylosaurus had 29 to 30 presacral vertebrae, 6 to 7 pygal vertebrae, and 89 to 112 caudal vertebrae; due to the lack of a bony articulation between the ilium and vertebral column, it is unclear whether any mosasaurs possessed true sacral vertebrae. In all tylosaurines, like in plioplatecarpines, the chevrons articulate to the caudal vertebrae, and are not fused to them, as they are in mosasaurines. The tail possesses a distinct downward curve, suggesting the presence of a tail fluke.
Soft tissue
Skin and coloration
Fossil evidence of the skin of Tylosaurus in the form of scales has been described since the late 1870s. These scales were small and diamond-shaped and were arranged in oblique rows, comparable to that found in modern rattlesnakes and other related reptiles. However, the scales in the mosasaur were much smaller in proportion to the whole body. An individual measuring in total body length had dermal scales measuring , and in each square inch (2.54 cm) of the mosasaur's underside an average of ninety scales were present. Each scale was keeled in a form resembling that of a shark's denticles. This probably helped reduce underwater drag and reflection on the skin.
Microscopic analysis of scales in a T. nepaeolicus specimen by Lindgren et al. (2014) detected high traces of the pigment eumelanin indicative of a dark coloration similar to the leatherback sea turtle in life. This may have been complemented with countershading, present in many aquatic animals, though the distribution of dark and light pigments in the species remains unknown. A dark-colored form would have provided several evolutionary advantages. Dark coloration increases absorption of heat, allowing the animal to maintain elevated body temperatures in colder environments. Possession of this trait during infancy would in turn facilitate fast growth rates. Unreflective dark coloring and countershading would have provided the mosasaur with increased camouflage. Additional speculative functions includes increased tolerance to solar ultraviolet radiation, strengthened integuments. The study remarked that certain melanism-coding genes are pleiotropic for increased aggression.
Respiratory system
AMNH FR 221 preserves parts of the cartilaginous respiratory system. This includes parts of the larynx (voice box), trachea (windpipe), and bronchi (lung airways). They were however only briefly described in the preserved position by Osborn (1899). The larynx is poorly preserved; a piece of its cartilage first appears below just between the pterygoid and quadrate and extends to behind the latter. This connects to the trachea, which appears below the atlas vertebra but is not preserved afterwards. The respiratory tract reappears below the fifth rib as a pair of bronchi and extends to just behind the as-preserved coracoids where preservation is lost. The pairing is suggestive of two functional lungs like modern limbed lizards but unlike snakes. Similar branching is also found in Platecarpus and putatively Mosasaurus, the only two other derived mosasaurs with their respiratory systems documented. The bifurcation point for the Tylosaurus specimen is anywhere between the first and sixth cervical vertebrae. In Platecarpus, the bronchi probably diverged below the sixth cervical into near-parallel pairs, while in Mosasaurus the organ is dislocated. A bifurcation point's position ahead of the forelimbs would be unlike terrestrial lizards, whose point is within the chest region, but similar to the short trachea and parallel bronchi of whales.
Classification
Taxonomy
Tylosaurus is classified within the family Mosasauridae in the superfamily Mosasauroidea. The genus is the type genus of its own subfamily, the Tylosaurinae. Other members of this group include Taniwhasaurus and possibly Kaikaifilu, and the subfamily is defined by a shared feature of an elongated premaxillary rostrum that does not bear teeth. The closest relatives of the Tylosaurinae include the Plioplatecarpinae and the primitive subfamilies Tethysaurinae and Yaguarasaurinae; together they are members of one of three possible major lineages of mosasaurs (the others being the Mosasaurinae subfamily and Halisauromorpha group) that was first recognized in 1993. This clade was named the Russellosaurina by Polcyn and Bell in 2005.
Tylosaurus was among the earliest derived mosasaurs. The oldest fossil attributable to the genus is a premaxilla (TMM 40092-27) recovered from Middle Turonian deposits of the Arcadia Park Shale in Texas, which is dated between 92.1 and 91.4 million years old based on correlations with index fossils. Although formally referred to as Tylosaurinae incertae sedis during its first description, it was remarked to probably belong to T. kansasensis. The specimen was later listed within the species in a 2020 reexamination. A slightly younger specimen is of a skull (SGM-M1) of an indeterminate Tylosaurus species similar to T. kansasensis from the Ojinaga Formation in Chihuahua, Mexico, dated around ~90 million years old at earliest. A tooth from a Late Maastrichtian deposit in Nasiłów, Poland dating close to the Cretaceous–Paleogene boundary has been attributed to Hainosaurus sp. With the incorporation of Hainosaurus as a synonym of Tylosaurus, this also makes the genus one of the last mosasaurs. Currently, eight species of Tylosaurus are recognized by scientists as taxonomically valid. They are as follow: T. proriger, T. nepaeolicus, T. bernardi, T. gaudryi, T. ivoensis, T. iembeensis, T. pembinensis, and T. saskatchewanensis. The validity of two additional taxa remain unsettled; there is still debate whether T. kansasensis is synonymous with T. nepaeolicus, and T. "borealis" has yet to be described in a formal publication.
Phylogeny and evolution
In 2020, Madzia and Cau performed a Bayesian analysis to better understand the evolutionary influence on early mosasaurs by contemporaneous pliosaurs and polycotylids by examining the rates of evolution in mosasauroids like Tylosaurus (specifically T. proriger, T. nepaeolicus, and T. bernardi). A Bayesian analysis in the study's implementation can approximate numerically defined rates of morphological evolution and ages of divergence of clades. The Tylosaurinae was approximated to have diverged from the Plioplatecarpinae around 93 million years ago; the divergence was characterized by the highest rate of evolution among all mosasaurid lineages. This trend of rapid evolution coincided with the extinction of the pliosaurs and a decrease in polycotylid diversity. The study noted converging traits between Tylosaurus, pliosaurs, and some polycotylids in tooth morphology and body size. However, there was no evidence to suggest that Tylosaurus or its precursors evolved as a result of out-competing and/or driving to extinction the pliosaurs and polycotylids. Instead, Madiza and Cau proposed that Tylosaurus may have taken advantage of the extinction of the pliosaurs and decline of polycotylids to quickly fill the ecological void they left behind. The Bayesian analysis also approximated a divergence of T. nepaeolicus from the rest of the genus around 86.88 million years ago and a divergence between T. proriger and T. bernardi around 83.16 million years ago. The analysis also generated a paraphyletic status of the genus, approximating Taniwhasaurus to have diverged from Tylosaurus around 84.65 million years ago, but this result is not consistent with previous phylogenetic analyses.
In the Western Interior Seaway, two species—T. nepaeolicus and T. proriger—may represent a chronospecies, in which they make up a single lineage that continuously evolves without branching in a process known as anagenesis. This is evident by how the two species do not stratigraphically overlap, are sister species, share minor and intermediate morphological differences such as a gradual change in the development of the quadrate bone, and lived in the same locations. The means by which this lineage evolved has been hypothesized to be through one of two evolutionary mechanisms related to changes in ontogeny. First, Jiménez-Huidobro, Simões, and Caldwell proposed in 2016 that T. proriger evolved as a paedomorph of T. nepaeolicus, in which the descendant arose as a result of morphological changes through the retention of juvenile features of the ancestor in adulthood. This was based on the presence of a frontal crest and convex borders of the parietal bone of the skull shared in both juvenile T. nepaeolicus and all T. proriger but lost in adult T. nepaeolicus. However, an ontogenetic study by Zietlow (2020) found that it was unclear whether this observation was a result of paedomorphosis, although this uncertainty may have been due that the sample size of mature T. nepaeolicus was too low to determine statistical significance. Second, the same study proposed an alternative hypothesis of peramorphosis, in which T. proriger evolved by developing traits found in mature T. nepaeolicus during immaturity. Based on results from a cladistical ontogram developed using data from 74 Tylosaurus specimens, the study identified a multitude of traits that were present in all T. proriger and mature T. nepaeolicus but absent in juvenile T. nepaeolicus: the skull size and depth are large, the length of the elongated rostrum exceeds 5% of the total skull length, the quadrate suprastapedial processes are thick, the overall quadrate shape converges, and the posteroventral process is fan-like.
The following cladogram is modified from a phylogenetic analysis by Jiménez-Huidobro & Caldwell (2019) using Tylosaurus species with sufficiently known material to model accurate relationships; T. gaudryi, T. ivoensis, and T. iembeensis were excluded from the analysis due to extensive missing data (i.e., lack of material with scoreable phylogenetic characters).
Paleobiology
Growth
Konishi and colleagues in 2018 assigned specimen FHSM VP-14845, a small juvenile with an estimated skull length of , to Tylosaurus based on the shape of the premaxilla, the proportions of the basisphenoid, and the arrangement of the teeth on the pterygoid. However, the specimen lacks the characteristically long premaxillary rostrum of other Tylosaurus, which is present in juveniles of T. nepaeolicus and T. proriger with skull lengths of . This suggests that Tylosaurus rostrum grew rapidly at an early stage in life, and also suggests that it did not develop due to sexual selection. Konishi and colleagues suggested a function in ramming prey, as employed by the modern orca.
Metabolism
Nearly all squamates are characterized by their cold-blooded ectothermic metabolism, but mosasaurs like Tylosaurus are unique in that they were likely endothermic, or warm-blooded. The only other known lizard with such a trait is the Argentine black and white tegu, though only partially. Endothermy in Tylosaurus was demonstrated in a 2016 study by Harrell, Pérez‐Huerta, and Suarez by examining δ18O isotopes in Tylosaurus bones. δ18O levels can be used to calculate the internal body temperature of animals, and by comparing such calculated temperatures between coexisting cold-blooded and warm-blooded animals, the type of metabolism can be inferred. The study used the body temperatures of the cold-blooded fish Enchodus and sea turtle Toxochelys (correlated with ocean temperatures) and warm-blooded seabird Ichthyornis from the Mooreville Chalk as a proxy. Analyzing the isotope levels of eleven Tylosaurus specimens an average internal body temperature of was calculated. This was much higher than the body temperature of Enchodus and Toxochelys ( and respectively) and similar to that of Ichthyornis (). Harrell, Pérez‐Huerta, and Suarez also calculated the body temperatures of Platecarpus and Clidastes with similar numbers, and respectively. The fact that the other mosasaurs were much smaller in size than Tylosaurus and yet maintained similar body temperatures made it unlikely that Tylosauruss body temperature was the result of another metabolic type like gigantothermy. Endothermy would have provided several advantages to Tylosaurus such as increased stamina for foraging larger areas and pursuing prey, the ability to access colder waters, and better adaptation to withstand the gradual cooling of global temperatures during the Late Cretaceous.
Mobility
Scientists previously interpreted Tylosaurus as an anguilliform swimmer that moved by undulating its entire body like a snake due to its close relationship with the animal. However, it is now understood that Tylosaurus actually used carangiform locomotion, meaning that the upper body was less flexible and movement was largely concentrated at the tail like in mackerels. A BS thesis by Jesse Carpenter published in 2017 examined the vertebral mobility of T. proriger spinal columns and found that the dorsal vertebrae were relatively rigid but the cervical, pygal, and caudal vertebrae were more liberal in movement, indicating flexibility in the neck, hip, and tail regions. This contrasted with more derived mosasaurs like Plotosaurus, whose vertebral column was stiff up to the hip. Interestingly, an examination of a juvenile T. proriger found that its cervical and dorsal vertebrae were much stiffer than those in adult specimens. This may have been an evolutionary adaptation among young individuals; a more rigid tail-based locomotion is associated with faster speed, and this would allow vulnerable juveniles to better escape predators or catch prey. Older individuals would see their spine grow in flexibility as predator evasion becomes less important for survival.
Tylosaurus likely specialized as an ambush predator. It was lightweight for a mosasaur of its size, having a morphological build designed to vastly reduce body mass and density. Its pectoral and pelvic girdles and paddles, which are associated with weight, are proportionally small. Its bones were highly cancellous and were likely filled with fat cells in life, which also increased buoyancy. It is unlikely that the latter trait was evolved in response to increasing body size as the similarly sized Mosasaurus hoffmannii lacked highly cancellous bone. These traits allowed Tylosaurus to be more conservative in its energy requirements, which is useful when traveling between ambush sites over large distances or through stealth. In addition, a reduced body density likely helped Tylosaurus to rapidly accelerate during an attack, assisted with the long and powerful tail of the mosasaur.
A 1988 study by Judith Massare attempted to calculate the sustained swimming speed, the speed at which the animal moves without tiring, of Tylosaurus through a series of mathematical models incorporating hydrodynamic characteristics and estimations of locomotive efficiency and metabolic costs. Using two T. proriger specimens, one long and the other , she calculated a consistent average maximum sustained swimming speed of . However, when testing whether the models represented an accurate framework, they were found to exaggerated. This was primarily because the variables accounting for drag may have been underestimated; estimation of drag coefficients for an extinct species can be difficult as it requires a hypothetical reconstruction of the morphological dimensions of the animal. Massare predicted that the actual sustained swimming speed of Tylosaurus was somewhere near half the calculated speed.
Feeding
One of the largest marine carnivores of its time, Tylosaurus was an apex predator that exploited the wide variety of marine fauna in its ecosystem. Stomach contents are well documented in the genus, which includes other mosasaurs, plesiosaurs, turtles, birds, bony fish, and sharks. Additional evidence from bite marks suggests the animal also preyed on giant squid and ammonites.
The enormous and varied appetite of Tylosaurus can be demonstrated in a 1987 find that identified fossils of a mosasaur measuring or longer, the diving bird Hesperornis, a Bananogmius fish, and possibly a shark all within the stomach of a single T. proriger skeleton (SDSM 10439) recovered from the Pierre Shale of South Dakota. Other records of stomach contents include a sea turtle in a T. bernardi-like species, a long Dolichorhynchops in another ( long) T. proriger, partially digested bones and scales of a Cimolichthys in a third T. proriger, partially digested vertebrae of a Clidastes in a fourth T. proriger, remains of three Platecarpus individuals in a T. nepaeolicus, and Plioplatecarpus bones in a T. saskatchewanensis. Puncture marks on fossils of ammonites, the carapace of a Protostega, and the gladius of an Enchoteuthis have been attributed to Tylosaurus.
Pasch and May (2001) reported bite marks from a dinosaur skeleton known as the Talkeetna Mountains Hadrosaur, which was found in marine strata of the Turonian-age Matanuska Formation in Alaska. The features of these marks were found to closely match that of the teeth of T. proriger. Because the fossil's locality was of marine deposits, the study reasoned that the dinosaur must have drifted offshore as a bloat-and-float carcass that was subsequently scavenged by the mosasaur. It was unlikely that the marks were a result of predation, as that would have led to a puncture, preventing the buildup of the bloating gases that allowed the corpse to drift out to sea in the first place. Garvey (2020) criticized the lack of conclusive evidence to support this hypothesis and ruled out T. proriger as a possible culprit, given that the species did not appear until the Santonian and is exclusive to the Western Interior Seaway. However, close relatives did maintain a presence nearby, evidenced by fragmentary fossils of an indeterminate tylosaurine from Turonian deposits in the Russian Chukotsky District.
Social behavior
The behavior of Tylosaurus towards each other may have been mostly aggressive, evidenced by fossils with injuries inflicted by another of their own kind. Such remains were frequently reported by fossil hunters during the late 19th and early 20th centuries, but few examples reside as specimens in scientific collections. Many of these fossils consist of healed bite marks and wounds that are concentrated around or near the head region, implying that there were the result of non-lethal interaction, but the motives of such contact remain speculative. In 1993, Rothschild and Martin noted that some modern lizards affectionately bite their mate's head during courtship, which can sometimes result in injuries. Alternatively, they also observed that some males lizards also employ head-biting as territorial behavior against rivals in a show of dominance by grappling the head to turn over the other on its back. It is possible that Tylosaurus behaved in similar ways.
Lingham-Soliar (1992) noted suggestions that use of the combat-oriented elongated rostrum of Tylosaurus was not exclusive to hunting and that it may have also been applied in sexual behavior through battles over female mates between males. However, he observed the elongated rostrum was invariably present in all individuals regardless of sex, and subsequent studies by Konishi et al. (2018) and Zietlow (2020) confirmed this pattern. This would imply that sexual selection was not a driver in its evolution and instead refined through sex-independent means.
At least one fatal instance of intraspecific combat among Tylosaurus is documented in the T. kansasensis holotype FHSM VP-2295, representing a long animal, which possesses numerous injuries that indicate it was killed by a larger Tylosaurus. The skull roof and surrounding areas exhibit signs of trauma in the form of four massive gouges, and the dentary contains at least seven puncture wounds and gouges. These pathologies are characteristic of bite marks from a larger Tylosaurus that measured around in length. The largest of the marks are about in length, matching the size of large mosasaur teeth, and they are positioned along two lines that converge close to 30°, matching the angle that each jaw converges towards in a mosasaur skull. In addition, FHSM VP-2295 suffered damage to its neck: the cervical vertebrae were found articulated at an unnatural angle of 40° relative to the long axis of the skull. The pattern of preservation makes it unlikely that the condition of the vertebrae was a result of disturbances by scavengers and instead indicates damage caused by a violently twisted neck during life. In a reconstructed scenario, the larger Tylosaurus would have first attacked at an angle slightly below the left side of the victim's head. This impact would cause the victim's skull to roll to the right side, allowing the aggressor to sink its teeth into the skull roof and right lower jaw, crushing the jaw and causing further breaks of nearby bones, such as the pterygoid, and the twisting of the jaw outwards, which would cause the quadrate to detach from its position and for the spinal cord to twist and sever at the skull's base, leading to a swift death.
Paleopathology
Examining 12 North American Tylosaurus skeletons and one T. bernardi skeleton, Rothschild and Martin (2005) identified evidence of avascular necrosis in every individual. For aquatic animals, this condition is often a result of decompression illness, which is caused when bone-damaging nitrogen bubbles build up in inhaled air that is decompressed either by frequent deep-diving trips or by intervals of repetitive diving and short breathing. The studied mosasaurs likely gained avascular necrosis through such behaviors, and given its invariable presence in Tylosaurus it is likely that deep or repetitive diving was a general behavioral trait of the genus. The study observed that between 3-15% of vertebrae in the spinal column of North American Tylosaurus and 16% of vertebrae in T. bernardi were affected by avascular necrosis. Carlsen (2017) posited that Tylosaurus gained avascular necrosis because it lacked the necessary adaptations for deep or repetitive diving, although noted that the genus had well-developed eardrums that could protect themselves from rapid changes in pressure.
Unnatural fusion of some vertebrae in the tail has been reported in some Tylosaurus skeletons. A variation of these fusions may concentrate near the end of the tail to form a single mass of multiple fused vertebrae called a "club tail." Rothschild and Everhart (2015) surveyed 23 North American Tylosaurus skeletons and one T. bernardi skeleton and found that five of the North American skeletons exhibited fused tail vertebrae. The condition was not found in T. bernardi, but this does not rule out its presence due to the low sample size. Vertebral fusion occurs when the bones remodel themselves after damage from trauma or disease. However, the cause of such events can vary between individuals and/or remain hypothetical. One juvenile specimen with the club tail condition was found with a shark tooth embedded in the fusion, which confirms that at least some cases were caused by infections inflicted by predator attacks. The majority of vertebral fusion cases in Tylosaurus were caused by bone infections, but some cases may have alternatively been caused by any type of joint disease such as arthritis. However, evidence of joint disease was rare in Tylosaurus when compared to mosasaurs such as Plioplatecarpus and Clidastes. Similar amassing of remodeled bone is also documented in bone fractures in other body parts. One T. kansasensis specimen possesses two fractured ribs that fully healed. Another T. proriger skull shows a fractured snout, probably caused by ramming into a hard object such as a rock. Presence of some healing indicates that the individual survived for some extended time before death. The injury in a snout region containing many nerve endings would have inflicted extreme pain.
| Biology and health sciences | Prehistoric squamates | Animals |
2469650 | https://en.wikipedia.org/wiki/Scelidosaurus | Scelidosaurus | Scelidosaurus (; with the intended meaning of "limb lizard", from Greek / meaning 'rib of beef' and sauros/ meaning 'lizard') is a genus of herbivorous armoured ornithischian dinosaur from the Jurassic of England.
Scelidosaurus lived during the Early Jurassic Period, during the Sinemurian to Pliensbachian stages around 191 million years ago, at the time when Europe formed an island archipelago. Its fossils have been found in the Charmouth Mudstone Formation near Charmouth in Dorset, England, and these fossils are known for their excellent preservation. Scelidosaurus has been claimed as one of the earliest complete dinosaur, and is among the most completely known dinosaur of the British Isles. Despite this, a modern description only materialised in 2020. After initial finds in the 1850s, comparative anatomist Richard Owen named and described Scelidosaurus in 1859. Only one species, Scelidosaurus harrisonii named by Owen in 1861, is considered valid today, although one other species was proposed in 1996.
Scelidosaurus was about long. It was a largely quadrupedal animal, feeding on low scrubby plants, the parts of which were bitten off by the small, elongated head to be processed in the large gut. Scelidosaurus was lightly armoured, protected by long horizontal rows of keeled oval scutes that stretched along the neck, back and tail.
One of the oldest known and most "primitive" of the thyreophorans, the exact placement of Scelidosaurus within this group has been the subject of debate for nearly 150 years. This was not helped by the limited additional knowledge about the early evolution of armoured dinosaurs. Today most evidence suggests that Scelidosaurus is the most derived of the known basal thyreophorans, either closely related to Ankylosauria or Stegosauria+Ankylosauria.
Description
Size and posture
A full-grown Scelidosaurus was rather small compared to most later non-avian dinosaurs, but it was a medium-sized species in the Early Jurassic. Some scientists have estimated a length of 4 metres (13 ft). In 2010, Gregory S. Paul gave a body length of 3.8 metres (12.5 ft) and a weight of . Scelidosaurus was quadrupedal, with the hindlimbs longer than the forelimbs. It may have reared up on its hind legs to browse on foliage from trees, but its arms were relatively long, indicating a mostly quadrupedal posture. A trackway from the Holy Cross Mountains of Poland shows a scelidosaur like animal walking in a bipedal manner, hinting that Scelidosaurus may have been more proficient at bipedalism than previously thought.
Distinguishing traits
The first modern diagnosis was provided by David Bruce Norman in 2020. In a first article, Norman provided autapomorphies, unique derived characters, of the skull. The front snout bones, the premaxillae, have a common central rough extension, in life bearing a small upper beak. The nasal bone has on its upper outside a facet touching the inner side of the ascending branch of the premaxilla. The antorbital fenestra is present as a bean-shaped depression, its lower edge formed by a sharp ridge. The central parietal crest on the skull roof is formed by two parallel crests separated by a narrow trough on the midline. The roof of the nasal cavity is formed by special plates above the vomers, called the "epivomers". The epipterygoid bone is shaped as a small conical vertical structure of which the base connects to the upper side of the pterygoid bone by means of a lateral flat surface. The basioccipital has large oblique facets on the lower sides. The opisthotic has an expanded pedicel with facets on its underside. Elongated epistyloid bones project obliquely to the rear and below, from the back of the skull. A small spur-like structure on the upper edge of the paroccipital process encases the posttemporal fenestra. The rear of the skull is fused on its upper edge with a pair of large curved horn-shaped osteoderms. The lower jaw shows only little exostosis, limited to the angular, and lacking an attached osteoderm.
Skull
The head of Scelidosaurus was small, about twenty centimetres long, and elongated. The skull was low in side view and triangular in top view, longer than it was wide, similar to that of earlier ornithischians. The snout, largely formed by the nasal bones, was flat on top. Scelidosaurus still had the five pairs of fenestrae (skull openings) seen in basal ornithischians: apart from the nostrils and eye sockets which are present in all basal dinosaurs, the fenestra antorbitalis and the upper and lower temporal fenestrae were not closed or overgrown, as with many later armoured forms. In fact, the upper temporal fenestrae were very large, forming conspicuous round openings in the top of the rear skull, serving as attachment areas for the powerful muscles that closed the lower jaws. The eye socket was slightly overshadowed in its front part by a brow ridge that has been seen as the prefrontal bone. In 2020, Norman concluded that it was a fused palpebral bone. Behind it, the upper rim of the eye socket was formed by the supraorbital bone. A study by Susannah Maidment e.a. concluded that juvenile specimens show that this bone was a fusion of three elements, one in front, the next in rear, and the third at the inner side.
The premaxilla, the bone forming the snout tip, was short and no predentary, the bone core of the lower beak on the tip of the stout lower jaws, has been found, so the horny beak that is assumed present with all ornithischians was likely very short. Its teeth were longer and more triangular in side view than in later armoured dinosaurs. There were at least five teeth in each premaxilla, and at least nineteen in the maxilla and sixteen in the dentary of the lower jaw. However, the number of maxillary and dentary teeth were established with the incomplete skull of one of the first specimens found; the actual numbers might have ranged up to about two dozen, perhaps twenty-six for the lower jaw. The premaxillary teeth were somewhat longer and recurved. To the rear, they gradually approach the form of the maxillary teeth, beginning to show denticles. The crowns of the maxillary and dentary teeth have denticles on their edges and a swollen basis
The ascending branches of the paired premaxillae notched the combined nasal bones, whereas the opposite was usual in ornithischians. The frontal bones were covered by a halo of fine ridges; these indicate the presence of keratinous plates, as with modern turtles. At the front of the braincase, paired hatchet-shaped ossified orbitosphenoids formed the floor of the olfactory lobes of the brain. The skull of the lectotype was damaged by a paleoichthyologist resulting in the detachment of triangular plates from the palate. These elements had been sketched by Norman in the seventies prior to the incident and interpreted as parts of the pterygoids, but in 2020 he concluded that they were special bones covering the roof of the nasal cavity, which he named the "epivomers". These are not known from any other animal.
Postcranial skeleton
The vertebral column of Scelidosaurus contained at least six neck vertebrae, seventeen dorsal vertebrae, four sacral vertebrae and at least thirty-five tail vertebrae.
Though perhaps the actual total of cervical vertebrae was as high as seven or eight, the neck was only moderately long. The torso was relatively flat in side view, however, despite the belly being broad, it was not extremely vertically compressed as with ankylosaurs but taller than wide. The last three dorsal vertebrae had no ribs. The spines of the sacral vertebrae touched each other but were not fused into a supraneural plate. The quickly tapering tail was relatively short, probably representing about half of body length. The tail chevrons were strongly inclined to the rear. The hip area and tail base were stiffened by large numbers of ossified tendons.
The scapula was short with a moderately expanded upper end. The coracoid was circular in side view. The elements of the forelimb were generally moderately long, straight and stout. The hand is only known from recent discoveries and has not yet been described. In the rather wide pelvis, the ilium was straight in side view. Its front blade was rod-shaped and moderately splayed to the outside, creating room for the belly. This was reinforced by the sacral ribs becoming longer towards the front. The sacral ribs were wider at their attachment areas with the ilium, but were not fused into a sacral yoke. The pubis featured a short prepubis. The pubis shaft was straight, running parallel to a straight ischium shaft that was transversely flattened at its lower end. The thighbone was straight in side view, in front view it was somewhat bowed to the outside. Its head was not separated from the shaft by a real neck. While the major trochanter was at about the same level as the head, the lower minor trochanter was separated from both by a deep cleft. At it rear side, the femur mid-shaft featured a well-developed drooping fourth trochanter, a process for the attachment of the retractor tail muscle, the Musculus caudofemoralis longus. The lower leg was somewhat shorter than the thighbone. The tibia had a wide upper end, with a cnemial crest protruding well to the front. The tibia lower end was also robust and rotated about 70° compared to the upper part, turning the foot strongly to the outside. The foot was very large and wide. The fifth metatarsal was only rudimentary but the other four were robust. Scelidosaurus had four large toes, with the innermost digit being the smallest. The fourth metatarsal was short but its toe was long and built to be splayed to the outside of the foot, to improve the stability. The claws were flat, hoof-shaped and curved to the inside.
Armour
The most obvious feature of Scelidosaurus is its armour, consisting of bony scutes embedded in the skin. These osteoderms were arranged in horizontal parallel rows down the animal's body. Osteoderms are today found in the skin of crocodiles, armadillos and some lizards. The osteoderms of Scelidosaurus ranged in both size and shape. Most were smaller or larger oval plates with a high keel on the outside, the highest point of the keel positioned more to the rear. Some scutes were small, flat and hollowed-out at the inside. The larger keeled scutes were aligned in regular horizontal rows. There were three rows of these along each side of the torso. The scutes of the lowest, lateral, row were more conical, rather than the blade-like osteoderms of Scutellosaurus. Between these main series, one or two rows of smaller oval keeled scutes were present. There were in total four rows of large scutes on the tail: one at the top midline, one at the midline of the underside, and one at each tail side. Whether the midline tail scutes continued over the torso and neck to the front is unknown and unlikely for the neck, though Scelidosaurus is often pictured this way.
The skull consists of bony outgrowths covering much of its surface, accompanied by separate bony plates encircling the upper portion of the eye socket and forming a pair of horns that extend upward and backward from the rear of the skull. The neck had at each side two rows of large scutes. The osteoderms of the lower neck row were very large, flat and plate-like. In this area the dermal armor includes a framework of deep foundational plates beneath more superficial bony elements. These foundational plates function as growth zones and anchors for their overlying components, thickening and spreading as they grow. The first osteoderms of the top neck rows formed a pair of unique three-pointed scutes directly behind the head. These points seem to have been connected by tendons to the rear joint processes, the postzygapophyses, of the axis vertebra.
Their edges develop grooves and eventually interlock, resulting in collar-like structures made of interconnected sections adorned with bony projections. These structures form four paired segments that shield the upper and side areas of the neck. Whether these segments fused along the centerline to create continuous neck rings, as observed in some related species, remains uncertain. The bony plates of the neck display diverse shapes, ranging from small and pointed to flat and cap-like, tall and ridged, or broad and blade-like. Behind the neck, the bony projections lack foundational plates and are arranged in three main rows running lengthwise along the body. The largest are positioned along the sides and are uniformly oval, ridged, and hollow. Smaller elements fill spaces between these primary rows, while the body’s remaining areas are covered by tightly packed, tiny bony structures, many of which supported keratinized scales. In general the scutes were larger at the front of the torso, the osteoderms diminishing towards the rear, especially on the surface of the thighs. The smallest flat round scutes might have filled the room between the larger osteoderm rows. Perhaps a row of vertical osteoderms was present on the upper arms. Compared to the later Ankylosauria, Scelidosaurus was lightly armoured, without continuous plating, spikes or pelvic shield. Rough areas on the skull and lower jaws indicate the presence of skin ossifications.
The limbs were sheathed in a mosaic of small bony nodules and rows of narrow, elongated structures with ridges. The tail exhibits a distinct arrangement, with four rows of tall, hollow, ridged bony elements along its top, bottom, and sides. Near the base of the tail, the underside features closely spaced, deep bony keels. The variation in bony armor across this animal, particularly in the neck and sides, suggests that larger projections were likely covered by keratin, while smaller, hollow forms may have served other functional or adaptive roles. Some of the latest specimens found show partly different osteoderms including scutes on which the keel is more like a thorn or spike. These specimens also seem to have little horns on the rear corners of the head, placed on the squamosal bones. Fossilized skin impressions have also been found. Between the bony scutes, Scelidosaurus had rounded non-overlapping scales like the present Gila monster. Between the large scutes, very small (5-10 millimetres [0.2-0.4 in]) flat "granules" of bone were perhaps distributed within the skin. In the later Ankylosauria, these small scutes may have developed into larger scutes, fusing into the multi-osteodermal plate armour seen in genera such as Ankylosaurus.
History of study
Discovery, naming, and type specimens
During the 1850s, quarry owner James Harrison of Charmouth, West Dorset of England found fossils from the cliffs of Black Ven between Charmouth and Lyme Regis, that were quarried, possibly for raw material for the manufacture of cement. Some of these he gave to the collector and retired general surgeon Henry Norris. In 1858, Norris and Harrison sent some fragmentary limb bones to Professor Richard Owen of the British Museum (Natural History), London (today the Natural History Museum). Among them was a left thighbone, specimen GSM 109560. In 1859, Owen named the genus Scelidosaurus in an entry about palaeontology in the Encyclopædia Britannica. The lemma text contained a diagnosis, implicating that the genus was validly named and was not a nomen nudum, despite the fact that the definition was vague and no specimens were identified. Owen intended to call the dinosaur "hindlimb saurian" but confused the Greek word σκέλος, , "hindlimb", with σκελίς, , "rib of beef". The name was inspired by the strong development of the hind leg. Afterwards Harrison sent a knee joint, a claw (GSM 109561), a juvenile specimen and a skull to Owen, that were described in 1861. On that occasion the type species Scelidosaurus harrisonii was named, the specific name honouring Harrison. The skull later was revealed to be part of a nearly complete skeleton, that was described by Owen in 1863.
British palaeontologist David Bruce Norman has stressed how remarkable it is that Owen, who previously had propounded that dinosaurs were active quadrupedal animals, largely neglected Scelidosaurus though it could serve as a prime example of this hypothesis and its fossil was one of the most complete dinosaurs found at that time. Norman explained this by Owen's excessive workload in this period, including several administrative functions, polemics with fellow-scientists and the study of a large number of even more interesting newly discovered extinct animals, such as Archaeopteryx. Norman also pointed out that Owen in 1861 suggested a lifestyle for Scelidosaurus that is very different from present ideas: it would have been a fish-eater and partially sea-dwelling.
Owen had not indicated a holotype. In 1888, Richard Lydekker while cataloguing the NHMUK fossils, designated some of the hindlimb fragments described in 1861, specimen NHMUK PV 39496 consisting of a lower part of a femur and an upper part of the tibia and fibula, together forming a knee joint, as the type specimen, hereby implicitly choosing them as the lectotype of Scelidosaurus. Lydekker gave no reason for this choice; perhaps he was motivated by their larger size. Unfortunately, mixed in with the Scelidosaurus fossils had been the partial remains of a theropod dinosaur and the femur and tibia thus belonged to such a carnivore; this was not discovered until 1968 by Bernard Newman. The same year, B. H. Newman suggested to have Lydekker's selection of the knee joint as the lectotype officially rescinded by the International Commission on Zoological Nomenclature, as the joint was in his opinion from a species related to Megalosaurus. Eventually, after Newman had already died, Alan Jack Charig actually filed a request in 1992. In 1994 the ICZN reacted positively, in Opinion 1788 deciding that the skull and skeleton, specimen NHMUK PV R.1111, would be the new lectotype of Scelidosaurus. The knee joint was in 1995 by Samuel Welles et al. informally assigned to a "Merosaurus", which name has not yet been validly published. It more likely belongs to some member of the Coelophysoidea or Neoceratosauria. It has also been established by Newman and confirmed by Roger Benson that the original left thighbone, GSM 109560, belonged to a theropod.
The new lectotype skeleton had been uncovered in the Black Ven Marl or Woodstone Nodule Bed, marine deposits of the Charmouth Mudstone Formation, dating from the late Sinemurian stage, about 191 million years ago. It consists of a rather complete skeleton with skull and lower jaws. Only the snout tip, the neck base, the forelimbs and the tail end are missing. Hundreds of osteoderms were found in connection with the skeleton, many more or less in their original position. From the 1960s onward, this fossil was further prepared by Ronald Croucher using acid baths to free the bones from the surrounding matrix, a method perfected for the Charmouth fossils. In 1992, Charig reported that only a single block had yet to be treated, but he died before the results could be published. Norman, who intended to complete this task, had revealed some new anatomical details in 2004. Apart from these, a modern description was largely lacking. In 2020, Norman published articles on the skull and the postcrania, also taking later finds into account. It transpired that the acid baths had, through leakages, severely deteriorated the condition of the bones, further mishandling leading to breakage and crumbling.
Additional specimens
Apart from the lectotype, other fossils are known of Scelidosaurus. In 1888 Lydekker catalogued a large number of single bones, largely limb elements, and osteoderms, that had been acquired by the NHMUK from the Norris collection. Owen in 1861 described a second, partial, skeleton of a juvenile animal, that later was added to the collection of Elizabeth Philpot and today is registered in the Lyme Regis Museum as specimen LYMPH 1997.37.4-10. As it was relatively large, Owen speculated, in the context of its presumed marine lifestyle, that Scelidosaurus might have been ovoviviparous. The short prepubis in this specimen convinced scientists that this process did not represent the main pubic body as some had thought, who had been unable to believe that the thin, backward-pointing, pubis with the Ornithischia was homologous to the forward-pointing much larger pubic bone in most reptilian groups.
In more recent times, new discoveries have been made at Charmouth, not through commercial quarrying but by the efforts of amateur palaeontologists. In 1968 a second partial juvenile skeleton was described, specimen NHMUK PV R6704, that had already been reported in 1959. Found by geologist James Frederick Jackson (1894–1966) of Charmouth, it is from a slightly younger layer, the Stonebarrow Marl Member dating to the early Pliensbachian, about 190 million years old. In 1985 Simon Barnsley, David Costain and Peter Langham excavated a partial skeleton including a very complete skull and skin impressions, which was sold to the Bristol Museum where it is registered as specimen BRSMG CE12785. Specimen CAMSMX.39256 is part of the collection of the Sedgwick Museum at Cambridge. Several specimens remain undescribed because they are in private collections. These include a 3.1 metres (ten feet) long skeleton found by David Sole in 2000, perhaps the most complete non-avian dinosaur exemplar ever discovered in the British Isles. All elements of the skeleton are now known. The finds by Sole differ from the lectotype in details of the armour and might represent a separate taxon or reflect sexual dimorphism. In 2020, Norman denied this.
Between the years 1980 and 2000, three fossils were discovered on a beach near The Gobbins in Northern Ireland by palaeontologist Roger Byrne. Exact geologic provenance is not reported for any of the specimens, but the very dark colouration of the specimens indicate (through means of comparison to marine fossils in other Northern Irish localities) they hail from Lias Group rocks, likely from either the Planorbis Zone or the Pre-planorbis Zone of the Waterloo Mudstone Formation. The specimens include BELUM K3998, a proximal femur fragment discovered in January 1980; BELUM K12493, the fragment of a tibia shaft discovered in April 1981; and BELUM K2015.1.54, a small pentagonal object discovered in 2000. Histologist Robin Reid recognized the first specimen as dinosaurian due to its bone texture and structure, and reported it in 1989, suspecting it belonged to Scelidosaurus or a similar animal. Byrne then recognized the tibia specimen as dinosaurian using similar identifiers; it was assumed, based on association, the two specimens came from the same animal. The pentagonal specimen was then assumed to be a scelidosaur osteoderm on the same logic.
These Irish specimens, alongside another discovered by fossil collector William Gray sometime in the late 19th or early 20th century, were formally studied by Michael J. Simms and colleagues and a study was published on them in the journal Proceedings of the Geologists' Association in December 2021. The assignment of the femoral fragment was upheld, with a clear ornithischian identity and with size and morphology specifically very similar to Scelidosaurus and unlike close relative Scutellosaurus. However, the tibia was reinterpreted as that of an indeterminate neotheropod, the pentagonal object as a mere piece of basalt resembling a fossil, and Grey's specimen as belonging to an ichthyosaur. The scelidosaur femur and theropod tibia are the only known remains of dinosaurs from Ireland, which has a poor Mesozoic fossil record entirely consisting of marine localities, and the scelidosaur specimen was the first ever reported from the island. However, in 2024, Satchell reidentified the proximal femur fragment (BELUM K3998) as an indeterminate dinosaur remain, not belonging to a Scelidosaurus or an ornithischian.
In 2000, David Martill et al. announced the preservation of soft tissue in a specimen referred to a cf. Scelidosaurus sp. (that is, material tentatively referred to the genus Scelidosaurus but not to any specific species). The fossil, with inventory number BRSMG CF2781, was in the early 1990s, in an already prepared state, discovered in the legacy of the late Professor John Challinor, who had used it to illustrate his lectures. Its provenance is unknown. It consists of a series of eight caudal vertebrae in a cut slab of carbonate mudstone, which was judged to date from the late Hettangian to Sinemurian stages. Parts of the fossil were preserved in such a way that an envelope of preserved soft tissue is visible around the vertebrae, and show the presence of an epidermal layer over the scutes. The authors concluded that the osteoderms of all basal armoured dinosaurs were covered in a tough, probably keratinous layer of skin.
Additional species
Scelidosaurus harrisonii, named and described by Owen, is currently the only recognized species, based on several nearly complete skeletons. A potential second species from the Sinemurian-age Lower Lufeng Formation, Scelidosaurus oehleri, was described by David Jay Simmons in 1965 under its own genus, Tatisaurus. In 1996 Spencer G. Lucas moved it to Scelidosaurus. Although the fossils are fragmentary, this reassessment has not been accepted, and S. oehleri is today once again recognized as Tatisaurus.
In 1989, scutes which were found in the Kayenta Formation (Glen Canyon Group) of northern Arizona, were by Kevin Padian referred to a Scelidosaurus sp., and used to determine that the age of the strata was around 199.6-196.5 million years ago, at a time when it was still thought that Scelidosaurus harrisonii dated to the early Sinemurian. These scutes established a geographic tie-in between Arizona's Glen Canyon and Europe, where fossils of Scelidosaurus had previously been discovered. Later scientists have rejected the assignment to Scelidosaurus, as the scutes are different in form. In 2014, Roman Ulansky named a new species, S. arizonenesis, based on these specimens. In 2016, Peter Malcolm Galton and Kenneth Carpenter identified it as a nomen dubium, instead once again placing the specimens as Thyreophora indet.
Classification and phylogeny
Scelidosaurus was placed in the Dinosauria by Owen in 1861. In 1868/1869 Edward Drinker Cope proposed a family Scelidosauridae in a double lecture but this was only published in December 1871; therefore it was Thomas Henry Huxley who validly named the Scelidosauridae in 1869. In the nineteenth century almost any armoured dinosaur then known has been considered a member of the Scelidosauridae. In the later twentieth century, the term was used for an assembly of "primitive" ornithischians close to the ancestry of ankylosaurs and stegosaurs, such as Scutellosaurus, Emausaurus, Lusitanosaurus and Tatisaurus. Today, paleontologists usually consider the Scelidosauridae paraphyletic, thus not forming a separate branch or clade; however, Benton (2004) lists the group as monophyletic. The family was resurrected by Chinese paleontologist Dong Zhiming in his 2001 description of Bienosaurus, a thyreophoran sharing close affinities with Scelidosaurus.
Scelidosaurus was an ornithischian. It was the oldest ornithischian known until the description of Geranosaurus in 1911. During the twentieth century, it has been classified at different times as an ankylosaur or stegosaur. Alfred von Zittel (1902), William Elgin Swinton (1934), and Robert Appleby et al. (1967) identified the genus as a stegosaurian, though this concept then encompassed all armoured forms. In a 1968 paper, Romer argued it was an ankylosaur. In 1977, Richard Thulborn of the University of Queensland attempted to reclassify Scelidosaurus as an ornithopod similar to Tenontosaurus or Iguanodon. Thulborn argued Scelidosaurus was a lightly built bipedal dinosaur adapted for running. Thulborn's 1977 theories on the genus have since been rejected.
This debate is still ongoing; at this time, Scelidosaurus is considered to be either more closely related to ankylosaurids than to stegosaurids and, by extension, a true ankylosaur, or basal to the ankylosaur-stegosaur split. The stegosaur classification has fallen out of favor, but is seen in older dinosaur books. Cladistic analyses have invariably recovered a basal position for Scelidosaurus, outside of the Eurypoda. In their comprehensive review on ornithischian phylogeny, André Fonseca and colleagues recoved Scelidosaurus as a basal thyreophoran outside of Eurypoda and defined Scelidosauridae in the PhyloCode as "the largest clade containing Scelidosaurus harrisonii, but not Ankylosaurus magniventris and Stegosaurus stenops".
The position of Scelidosaurus according to a cladistic study of 2011 is shown by this cladogram:
In the 2022 monograph on Scelidosaurus by David Norman, a different relationship amongst thyreophorans was found, with Stegosauria being the most basal group, and Scelidosaurus being most closely related to Ankylosauria.
Fossil records of thyreophorans more basal than Scelidosaurus are sparse. The more "primitive" Scutellosaurus, also found in Arizona, was an earlier genus which was facultatively bipedal. A trackway of a possible early armoured dinosaur, from around 195 million years ago, has been found in France. Ancestors of these basal thyreophorans evolved from early ornithischians similar to Lesothosaurus during the Late Triassic.
Paleobiology
Diet
Like most other thyreophorans, Scelidosaurus is known to be herbivorous. However, while some later ornithischian groups possessed teeth capable of grinding plant material, Scelidosaurus had smaller, less complex leaf-shaped teeth suitable for cropping vegetation and jaws capable of only vertical movement, due to a short jaw joint. Paul Barrett concluded that Scelidosaurus fed with a puncture-crush system of tooth-on-tooth action, with a precise but simple up-and-down jaw movement, in which the food was mashed between the inner side of the upper teeth and the outer side of the lower teeth, without the teeth actually touching each other as shown by very long vertical wear facets on the lower teeth alone.
In this aspect, it resembled the stegosaurids, which also bore primitive teeth and simple jaws. Its diet would have consisted of ferns or conifers, as grasses did not evolve until late into the Cretaceous Period, after Scelidosaurus was long extinct.
Another similarity with the stegosaurs is the narrow head, which might indicate a selective diet consisting of high-quality fodder. However, Barrett pointed out that for an animal the size of Scelidosaurus, with a large gut allowing efficient fermentation, the intake of easily digestible food of high energetic value was less important than with smaller animals, that are often critically dependent on it. Norman concluded that Scelidosaurus fed on low scrubby vegetation, with a height up to one metre. Raising itself on its hindlimbs alone, could have vertically increased the feeding envelope and was perhaps anatomically possible, but Norman doubted it was a relevant part of its behaviour.
Palaeoenvironment
During the Early Jurassic, Europe formed an island archipelago. The Charmouth outcrop is adjacent to a number of tectonic highs which at this point in time formed islands, including the Welsh High (comprising what is now much of Wales) to the north, as well as the Cornubian Massif (including what is now Devon and Cornwall) to the west. The presence of Classopolis pollen in the sediments suggests that conifers belonging to the extinct family Cheirolepidiaceae grew on the islands.
| Biology and health sciences | Ornitischians | Animals |
2469652 | https://en.wikipedia.org/wiki/Dimorphodon | Dimorphodon | Dimorphodon ( ) was a genus of medium-sized pterosaur from Europe during the early Jurassic Period (about 201-191 million years ago). It was named by paleontologist Richard Owen in 1859. Dimorphodon means "two-form tooth", derived from the Greek () meaning 'two', () meaning 'shape' and () meaning 'tooth', referring to the fact that it had two distinct types of teeth in its jaws – which is comparatively rare among reptiles. The diet of Dimorphodon has been questioned among researchers, with earlier interpretations depicting it as an insectivore or a piscivore. Recent studies have suggested that Dimorphodon likely hunted small vertebrates, though it still would have consumed soft invertebrates like insects.
Description
The body structure of Dimorphodon displays many "primitive" characteristics, such as, according to Owen, a very small brain-pan and proportionally short wings. The first phalanx in its flight finger is only slightly longer than its lower arm. The neck was short but strong and flexible and may have had a membranous pouch on the underside. The vertebrae had pneumatic foramina, openings through which the air sacs could reach the hollow interior. Dimorphodon had an adult body length of long, with a 1.45 metre (4.6 ft) wingspan. The tail of Dimorphodon was long and consisted of thirty vertebrae. The first five or six were short and flexible, but the remainder gradually increased in length and were stiffened by elongated vertebral processes. The terminal end of the tail may have borne a Rhamphorhynchus-like tail vane, although no impressions have yet been found in Dimorphodon fossils to confirm this speculation.
Skull
Dimorphodon had a large, bulky skull approximately in length, whose weight was reduced by large openings separated from each other by thin bony partitions. Its structure, reminiscent of the supporting arches of a bridge, prompted Richard Owen to declare that, as far as achieving great strength from lightweight materials was concerned, no vertebra was more economically constructed; Owen saw the vertebrate skull as a combination of four vertebrae modified from the ideal type of the vertebra. The front of the upper jaw had four or five fang-like teeth followed by an indeterminate number of smaller teeth; the maxilla of all exemplars is damaged at the back. The lower jaw had five longer teeth and thirty to forty tiny, flattened pointed teeth, shaped like lancets. Many depictions give it a speculative puffin-like 'beak' because of similarities between the two animals' skulls.
History of discovery
The first fossil remains now attributed to Dimorphodon were found in England by fossil collector Mary Anning, at Lyme Regis in Dorset, UK in December 1828. This region of Britain is now a World Heritage Site, dubbed the Jurassic Coast; in it layers of the Blue Lias are exposed, dating from the Hettangian-Sinemurian. The specimen was acquired by William Buckland and reported in a meeting of the Geological Society on 5 February 1829. In 1835, after a thorough study by William Clift and William John Broderip, this report, strongly expanded, was published in the Transactions of the Geological Society, describing and naming the fossil as a new species. As was the case with most early pterosaur finds, Buckland classified the remains in the genus Pterodactylus, coining the new species Pterodactylus macronyx. The specific name is derived from Greek makros, "large" and onyx, "claw", in reference to the large claws of the hand. The specimen, presently NHMUK PV R1034, consisted of a partial and disarticulated skeleton on a slab, lacking the skull. Buckland in 1835 also assigned a piece of jaw from the collection of Elizabeth Philpot to P. macronyx. Later, the many putative species assigned to Pterodactylus had become so anatomically diverse that they began to be broken into separate genera.
In 1858, Richard Owen reported finding two new specimens, NHMUK PV OR 41212 and NHMUK PV R1035, again partial skeletons but this time including the skulls. Having found the skull to be very different from that of Pterodactylus, Owen assigned Pterodactylus macronyx its own genus, which he named Dimorphodon. His first report contained no description and the name remained a nomen nudum. In 1859, however, a subsequent publication by Owen provided a description. After several studies highlighting aspects of Dimorphodons anatomy, Owen in 1874 made NHMUK PV R1034 the holotype.
Meanwhile, though Dimorphodon is not a very common fossil, other fragmentary specimens were found. Some of these were acquired by Othniel Charles Marsh between 1873 and 1881 from the London fossil dealer Bryce McMurdo Wright. One of these had been recovered from early Jurassic strata at the south bank of the Severn river, at the Aust Cliff.
An additional species of Dimorphodon, D. weintraubi, was named by James Clark et al in 1998 from a partial skeleton recovered in siltstones from the site Huizachal Canyon in La Boca Formation in Tamaulipas, Mexico, from the Early Jurassic (Pliensbachian), where remains of sphenodontians, dinosaurs and mammaliaforms have also been found. It is known from the type specimen, IGM 3494 (Instituto Geológico de México, of the Universidad Nacional Autónoma de México), that comprises articulated pieces of the skeleton including the posterior part of skull, four cervical vertebrae, the scapulocoracoids, left humerus, partial right wing and right leg distal to mid tibiotarsus. This specimen is larger than D. macronyx and the well preserved foot of it shows that pterosaurs do not have a digitigrade posture in their hindlimbs, but that it have a plantigrade gait, as has been inferred from footprints. The name of the species is a homage to Dr. Robert L. Weintraub. Later studies considered this species not closely related to Dimorphodon macronyx, but an early relative of Anurognathidae.
Classification
In 1870, Seeley assigned Dimorphodon to its own family, Dimorphodontidae, with Dimorphodon as the only member. It was suggested in 1991 by the German paleontologist Peter Wellnhofer that Dimorphodon might be descended from the earlier European pterosaur Peteinosaurus. Later exact cladistic analyses are not in agreement. According to Unwin, Dimorphodon was related to, though probably not a descendant of, Peteinosaurus, both forming the clade Dimorphodontidae, the most basal group of the Macronychoptera and within it the sister group of the Caelidracones. This would mean that both dimorphodontid species would be the most basal pterosaurs known with the exception of Preondactylus. According to Alexander Kellner, however, Dimorphodon is far less basal and not a close relative of Peteinosaurus.
The cladogram recovered by Andres and Myers in 2013 is reproduced below.
Palaeobiology
Diet
The knowledge of how Dimorphodon lived is limited. It perhaps mainly inhabited coastal regions and might have had a very varied diet. Buckland suggested it ate insects. Later, it became common to depict it as a piscivore (fish eater), though biomechanical studies support Buckland's original insectivore idea better, and inconsistent with the animal's habits (see flight below). Dimorphodon had an advanced jaw musculature specialized for a "snap and hold" method of feeding. The jaw could close extremely quickly, but with relatively little force or tooth penetration. This, along with the short and high skull and longer, pointed front teeth suggest that Dimorphodon was an insectivore, though it may have occasionally eaten small vertebrates and carrion as well. Mark Witton has argued that the animal was a specialised carnivore that was too large for an insectivorous diet, though he did acknowledge that it still might have ate large insects, and thus specialised to hunt relatively small vertebrates, with its relatively weak jaw musculature indicating that it probably ate proportionally small prey. Dental microwear examinations confirm its status as a vertebrate predator, as opposed to several other insectivore or piscivore early pterosaurs, though the study does acknowledge that the possibility of consuming relatively softer invertebrates cannot be excluded entirely.
Locomotion
Like many pterosaurs, Dimorphodon has been perceived as a soarer in the past, correlating to historical perceptions of pterosaurs as seabird analogues. However, more recent studies show that the animal was actually a rather poor flyer: its wings are proportionally short in relation to the body and its skeleton rather robust, offering very little gliding potential. In life, Dimorphodon probably relied on frantic short flights in the same manner as modern fowl, tinamous and woodpeckers, being unable to fly for long distances and probably only taking to the air as a last resort.
Its derived position amidst primitive pterosaurs implies that this ineptitude is a developed trait, not an ancestral characteristic, as earlier pterosaurs like Preondactylus were capable aeronauts.
Owen saw Dimorphodon as a quadruped. He speculated that the fifth toe supported a membrane between the tail and the legs and that the animal was therefore very ungainly on the ground. However, his rival Harry Govier Seeley, propagating the view that pterosaurs were warm-blooded and active, argued that Dimorphodon was either an agile quadruped or even a running biped due to its relatively well developed hindlimbs and characteristics of its pelvis. This hypothesis was revived by Kevin Padian in 1983. However, fossilised track remains of other pterosaurs (ichnites) show a quadrupedal gait while on the ground and these traces are all attributed to derived pterosaurs with a short fifth toe. Dimorphodon's was elongated, clawless, and oriented to the side. David Unwin has therefore argued that even Dimorphodon was a quadruped, a view confirmed by computer modelling by Sarah Sangster.
Like most non-pterodactyloid pterosaurs, Dimorphodon was a competent climber, possessing proportionally large and curved ungals and a low center of gravity. Like modern squirrels, it probably moved in a saltatorial manner as it climbed.
| Biology and health sciences | Pterosaurs | Animals |
2470309 | https://en.wikipedia.org/wiki/Nurek%20Dam | Nurek Dam | The Nurek Dam (; Tajik: Нерӯгоҳи обии Норак, Nerūgohi obii Norak, Tajik for Nurek Hydro-electric Station) is an earth-fill embankment dam on the Vakhsh River in Tajikistan. Its primary purpose is hydroelectric power generation and its power station has an installed capacity of 3,015 MW. Construction of the dam began in 1961 and the power station's first generator was commissioned in 1972. The last generator was commissioned in 1979 and the entire project was completed in 1980 when Tajikistan was still a republic within the Soviet Union, becoming the tallest dam in the world at the time. At , it is currently the second tallest man-made dam in the world, after being surpassed by Jinping-I Dam in 2013. The Rogun Dam, also along the Vakhsh in Tajikistan, may exceed it in size when completed.
Construction
The Nurek Dam was constructed by the Soviet Union between the years 1961 and 1980. It is uniquely constructed, with a central core of cement forming an impermeable barrier within a -high rock and earth fill construction. The volume of the mound is 54 million m3. The dam includes nine hydroelectric generating units, the first commissioned in 1972 and the last in 1979. An estimated 5,000 people were resettled from the dam's flooding area.
The dam is located in a deep gorge along the Vakhsh River in western Tajikistan, about east of the nation's capital of Dushanbe. A town near the dam, also called Nurek, houses engineers and other workers employed at the dam's power plant.
Electricity generation
A total of nine Francis turbine-generators are installed in the Nurek Dam's power station. Originally having a generating capacity of 300 MW each (2,700 MW total), they were redesigned and retrofitted between 1984 and 1988 so now have a capacity of 335 MW each (3,015 MW total). As of 1994, this formed most of the nation's 4.0 gigawatt hydroelectric generating capacity, which was adequate to meet 98% of the nation's electricity needs. As of early 2024, it supplies 70% of Tajikistan's electricity.
Reservoir
The reservoir formed by the Nurek Dam, known simply as Nurek, is the largest reservoir in Tajikistan with a capacity of 10.5 km3. The reservoir is over in length, and has a surface area of . The reservoir drives the hydroelectric plant located within the dam. Stored water is also used for irrigation of agricultural land. Irrigation water is transported 14 kilometers through the Dangara irrigation tunnel and is used to irrigate about of farmland. It is suspected that the reservoir may have caused induced seismicity when being impounded.
| Technology | Dams | null |
2472666 | https://en.wikipedia.org/wiki/Interplanetary%20dust%20cloud | Interplanetary dust cloud | The interplanetary dust cloud, or zodiacal cloud (as the source of the zodiacal light), consists of cosmic dust (small particles floating in outer space) that pervades the space between planets within planetary systems, such as the Solar System. This system of particles has been studied for many years in order to understand its nature, origin, and relationship to larger bodies. There are several methods to obtain space dust measurement.
In the Solar System, interplanetary dust particles have a role in scattering sunlight and in emitting thermal radiation, which is the most prominent feature of the night sky's radiation, with wavelengths ranging 5–50 μm. The particle sizes of grains characterizing the infrared emission near Earth's orbit typically range 10–100 μm. Microscopic impact craters on lunar rocks returned by the Apollo Program revealed the size distribution of cosmic dust particles bombarding the lunar surface. The ’’Grün’’ distribution of interplanetary dust at 1 AU, describes the flux of cosmic dust from nm to mm sizes at 1 AU.
The total mass of the interplanetary dust cloud is approximately , or the mass of an asteroid of radius 15 km (with density of about 2.5 g/cm3). Straddling the zodiac along the ecliptic, this dust cloud is visible as the zodiacal light in a moonless and naturally dark sky and is best seen sunward during astronomical twilight.
The Pioneer spacecraft observations in the 1970s linked the zodiacal light with the interplanetary dust cloud in the Solar System. Also, the VBSDC instrument on the New Horizons probe was designed to detect impacts of the dust from the zodiacal cloud in the Solar System.
Origin
The sources of interplanetary dust particles (IDPs) include at least: asteroid collisions, cometary activity and collisions in the inner Solar System, Kuiper belt collisions, and interstellar medium grains (Backman, D., 1997). The origins of the zodiacal cloud have long been subject to one of the most heated controversies in the field of astronomy.
It was believed that IDPs had originated from comets or asteroids whose particles had dispersed throughout the extent of the cloud. However, further observations have suggested that Mars dust storms may be responsible for the zodiacal cloud's formation.
Life cycle of a particle
The main physical processes "affecting" (destruction or expulsion mechanisms) interplanetary dust particles are: expulsion by radiation pressure, inward Poynting-Robertson (PR) radiation drag, solar wind pressure (with significant electromagnetic effects), sublimation, mutual collisions, and the dynamical effects of planets (Backman, D., 1997).
The lifetimes of these dust particles are very short compared to the lifetime of the Solar System. If one finds grains around a star that is older than about 10,000,000 years, then the grains must have been from recently released fragments of larger objects, i.e. they cannot be leftover grains from the protoplanetary disk (Backman, private communication). Therefore, the grains would be "later-generation" dust. The zodiacal dust in the Solar System is 99.9% later-generation dust and 0.1% intruding interstellar medium dust. All primordial grains from the Solar System's formation were removed long ago.
Particles which are affected primarily by radiation pressure are known as "beta meteoroids". They are generally less than 1.4 × 10−12 g and are pushed outward from the Sun into interstellar space.
Cloud structures
The interplanetary dust cloud has a complex structure (Reach, W., 1997). Apart from a background density, this includes:
At least 8 dust trails—their source is thought to be short-period comets.
A number of dust bands, the sources of which are thought to be asteroid families in the main asteroid belt. The three strongest bands arise from the Themis family, the Koronis family, and the Eos family. Other source families include the Maria, Eunomia, and possibly the Vesta and/or Hygiea families (Reach et al. 1996).
At least 2 resonant dust rings are known (for example, the Earth-resonant dust ring, although every planet in the Solar System is thought to have a resonant ring with a "wake") (Jackson and Zook, 1988, 1992) (Dermott, S.F. et al., 1994, 1997)
Rings of dust
Interplanetary dust has been found to form rings of dust in the orbital space of Mercury and Venus. Venus's orbital dust ring is suspected to originate either from yet undetected Venus trailing asteroids, interplanetary dust migrating in waves from orbital space to orbital space, or from the remains of the Solar System's circumstellar disc, out of which its proto-planetary disc and then itself, the Solar planetary system, formed.
Dust collection on Earth
In 1951, Fred Whipple predicted that micrometeorites smaller than 100 micrometers in diameter might be decelerated on impact with the Earth's upper atmosphere without melting. The modern era of laboratory study of these particles began with the stratospheric collection flights of Donald E. Brownlee and collaborators in the 1970s using balloons and then U-2 aircraft.
Although some of the particles found were similar to the material in present-day meteorite collections, the nanoporous nature and unequilibrated cosmic-average composition of other particles suggested that they began as fine-grained aggregates of nonvolatile building blocks and cometary ice. The interplanetary nature of these particles was later verified by noble gas and solar flare track observations.
In that context a program for atmospheric collection and curation of these particles was developed at Johnson Space Center in Texas. This stratospheric micrometeorite collection, along with presolar grains from meteorites, are unique sources of extraterrestrial material (not to mention being small astronomical objects in their own right) available for study in laboratories today.
Experiments
Spacecraft that have carried dust detectors include Helios, Pioneer 10, Pioneer 11, Ulysses (heliocentric orbit out to the distance of Jupiter), Galileo (Jupiter Orbiter), Cassini (Saturn orbiter), and New Horizons (see Venetia Burney Student Dust Counter).
Obscuring effect
The Solar interplanetary dust cloud obscures the extragalactic background light, making observations of it from the Inner Solar System very limited.
Major Review Collections
Collections of review articles on various aspects of interplanetary dust and related fields appeared in the following books:
In 1978 Tony McDonnell edited the book Cosmic Dust which contained chapters on comets along with zodiacal light as indicator of interplanetary dust, meteors, interstellar dust, microparticle studies by sampling techniques, and microparticle studies by space instrumentation. Attention is also given to lunar and planetary impact erosion, aspects of particle dynamics, and acceleration techniques and high-velocity impact processes employed for the laboratory simulation of effects produced by micrometeoroids.
2001 Eberhard Grün, Bo Gustafson, Stan Dermott, and Hugo Fechtig published the book Interplanetary Dust. Topics covered are: historical perspectives; cometary dust; near-Earth environment; meteoroids and meteors; properties of interplanetary dust, information from collected samples; in situ measurements of cosmic dust; numerical modeling of the Zodiacal Cloud structure; synthesis of observations; instrumentation; physical processes; optical properties of interplanetary dust; orbital evolution of interplanetary dust; circumplanetary dust, observations and simple physics; interstellar dust and circumstellar dust disks.
2019 Rafael Rodrigo, Jürgen Blum, Hsiang-Wen Hsu, Detlef V. Koschny, Anny-Chantal Levasseur-Regourd, Jesús Martín-Pintado, Veerle J. Sterken, and Andrew Westphal collected reviews in the book Cosmic Dust from the Laboratory to the Stars. Included are discussions of dust in various environments: from planetary atmospheres and airless bodies over interplanetary dust, meteoroids, comet dust and emissions from active moons to interstellar dust and protoplanetary disks. Diverse research techniques and results, including in-situ measurement, remote observation, laboratory experiments and modelling, and analysis of returned samples are discussed.
| Physical sciences | Solar System | Astronomy |
6018260 | https://en.wikipedia.org/wiki/Motorized%20scooter | Motorized scooter | A motorized scooter is a stand-up scooter powered by either a small internal combustion engine or electric hub motor in its front and/or rear wheel. Classified as a form of micromobility, they are generally designed with a large center deck on which the rider stands. The first motorized scooter was manufactured by Autoped in 1915.
Recently, electric kick scooters (e-scooters) have grown in popularity with the introduction of scooter-sharing systems that use apps to allow users to rent them by the minute; such systems were initially found in the United States and in Queensland, Australia, but now are in major cities and in all the western world.
History
1915: Autoped introduces its stand-up scooter. Pulling back on the handlebar disengaged the clutch and applied the brake. Production continued until 1921; Krupp of Germany built the Autoped under license from 1919 to 1922.
1986: Go-Ped introduces the first modern stand-up scooters, the Roadster and Sport.
May 2001: Go-Ped introduces the first full-suspension stand-up e-scooter, the Hoverboard.
2004: Evo Powerboards introduces the 2x, the first scooter with a two-speed transmission.
November 2009: Go-Ped introduces its first completely propane-powered scooter and go-kart, the GSR Pro-Ped and GSR Pro-Quad.
2009: Italian-Israeli designer Nimrod Ricardo Sapir designs the world's first folding e-scooter based on his patent.
2010: Nimrod Ricardo Sapir starts producing the world's first motorized folding e-scooter utilizing lithium-ion batteries and a brushless hub motor under the MyWay brand in Avihayil, Israel, renamed Inokim in 2013 and later moving production to Ningbo, China.
2013: Light electric folding scooters powered by rechargeable lithium batteries and brushless hub motors become available from Micro Mobility Systems AG.
2018: Dockless scooter-sharing systems are rolled out in major cities, largely as expansions of bike-sharing systems.
E-scooters
Overview
Usage
Motorized kick scooters are used in law enforcement, security patrolling and leisure. New ride-sharing systems have made e-scooters easily accessible. They are popular in urban areas and are used as an alternative to bicycling or walking. Ride sharing companies first started dropping these scooters off in large US cities in 2018, and the need for short distance easy access transportation in many cities has meant that they have become increasingly popular with more and more companies looking to join the market.
Environment
E-scooters, and other electric vehicles, have the potential to reduce carbon dioxide (CO2) emissions which are a cause of global warming, and other pollutants, if they are used to replace travel in vehicles with internal combustion engines. Potential environmental benefits depend upon how scooters are used: if they replace car journeys they may be beneficial, but not if they replace walked or cycled journeys. Manufacture of the batteries, in particular, requires resources, and they are often not recycled. Lime estimated that globally one in four trips on its scooters replaced a car journey. A December 2021 Swiss research paper found that privately owned e-scooters tended to replace car journeys, but rented e-scooters emitted more CO2 than the transport modes they replaced.
Safety
E-scooters are a potentially environmentally friendly alternative personal mode of transportation that has appeal in urban settings and for short distances. However, they are not exempt from the vulnerabilities users may encounter in road traffic injuries similar to exposures pedestrians and bicyclists have shared the roads. For example, Israel has seen over 120,000 imports of e-bike and e-scooters over a two-year period, but due to poor cycling infrastructure, cyclists are often forced onto pedestrian sidewalks, and pedestrians use bike lanes and thus increase the risk of traffic collision. A 2022 review of medical notes found that injury rates due to e-scooters were more like those of motorcycles than bicycles.
As availability and demand for e-scooters increases, with more powerful versions capable of reaching up to 50 miles per hour, the number of traffic accident cases has increased. Israel witnessed a six-fold increase of e-bike and e-scooter accidents over a span of three years, and China found a four-fold increase in injury rate and a six-fold increase in mortality rates. However, significant gaps remain in the knowledge about the safety measures and impact of e-scooters. A particular cause of accidents is the instability of vehicles with such small wheels when, for example, hitting a pothole.
As e-scooters become more popular in urban and high traffic settings, user safety poses a major concern alongside other health risks for drivers, pedestrians, cyclists and other vulnerable groups such as the elderly and children sharing the road. A study conducted in China assessed risky behaviors of e-bike, e-scooter, and bicycle riders at crossing signalized intersections and found three different types of risky behaviors including stopping beyond the stop line, riding in motor lanes, and riding against traffic. A study of 2014-2020 UCLA-affiliated hospitals and outpatient center visits found that e-scooter injury rates in Greater Los Angeles area were similar to those of motorcycles, with about 33% of victims needing extensive follow-up care. However, the fatality rate was comparable to pedal bikes.
The same study found that those riding e-scooters are more likely to engage in risky behaviors. In specific, e-scooter riders were more likely to ride in motor lanes and ride against the flow of traffic through there is high variability in the types of accidents that occur and can vary based on time of day. Under-reporting poses as additional gaps in knowledge, as minor crashes, for example, tend to be under-reported and thus unaccounted for in overall e-scooter injury prevalence and there exist gaps in research on injuries related to e-scooters. Scooter-sharing systems such as Lime or Bird include safety precautions on the scooters themselves, such as: "helmet required, license required, no riding on sidewalks, no double riding, 18+ years old". Apps used to unlock and rent the scooters will also have safety reminders and ask the riders to abide by local laws while using them. However, these recommendations are not always followed, and the difference in laws between cities and states makes regulation difficult.
A consumer association in Belgium tested e-scooters, concluding that a bicycle was preferable, citing many problems with the devices, including in particular battery failure and very poor braking in wet conditions. E-scooters were regulated as toys, without the safety considerations required for vehicles.
When electric kick scooters were introduced in Norway, the media reported a high increase in accidents, including several deaths.
In Britain as of late 2021 privately owned e-scooters could not be used on public roads or foot-ways; during a trial from mid-2020 until late 2022 rental scooters could be used on roads, but not foot-ways, by users with an appropriate driving licence. At the time private scooters were widely used, illegally, on foot ways and roads. There were safety concerns—scooter accidents were causing injuries more like motorcycles than pedal cycles. Privately owned scooters were banned from carriage on London public transport after a spate of battery fires.
Regulation
Australia
In Queensland, the laws around the use of e-scooters and other personal mobility devices are made and enforced by the state government.
While some local governments in Queensland have not allowed Lime Scooter trials, Brisbane City Council is currently undertaking a Lime Scooter trial and has invited tenders for two scooter contracts in the city.
In the ACT, the framework for personal mobility devices was amended to include e-scooters and other similar devices from 20 December 2019, permitting use on footpaths, shared paths, bicycle paths and the bicycle side of separated paths. Bicycle helmets are required to be worn.
Perth became the latest City to announce an e scooter trial, which launched in March 2023.
Austria
Electric vehicles with a power up to 600 watts and a speed up to 25 km/h are considered as bicycles.
Belgium
Belgium's traffic rules were updated on 1 June 2019 to be in line with the European Commission guidelines formed in 2016. It became legal for people over 15 years of age to ride electric motorized scooters with speed limited to 25 km/h on public roads, mirroring e-bikes. Protective gear and insurance are not required by law.
Canada
Commuting in Canada with an e-scooter has increased. As power-assisted bicycles, e-scooters must follow many of the same federal laws and regulations, such as being limited to 32 km/h and not being allowed over 500 W output. Ontario has recently unveiled a series of laws aimed at ensuring safety while using electric-kick scooters or, e-scooters. The new laws require all riders to carry a valid driver's license, and those under the age of 16 must be accompanied by an adult who also carries a valid driver's license. Riders are now also required to wear an approved helmet when operating their e-scooter and have bright lights installed on the front and back of their vehicles.
Denmark
Since 1 January 2022, helmets are mandatory.
Finland
In Finland e-scooters must comply with the same rules as bicycles and they do not have any age restrictions. However, all e-scooters that have a maximum speed over 25 km/h are classified as small motorcycles and require a motor insurance.
France
Currently France only allows e-scooters on footpaths if they have a maximum speed of . Those travelling at up to 25 km/h are relegated to bike lanes. Legislators are considering a new law that would force users of e-scooters going faster than to have a type A1 license—the same as for small motorcycles. The legal framework is very blurry and does not define where e-scooters may or may not be driven or parked. The Deputy Mayor of Paris Christophe Najdovski is lobbying Transport Minister Élisabeth Borne for a clearer framework that would give municipalities the power to tighten the rules on how permits are issued and how authorizations are given to deploy a fleet of e-scooters to operators.
French daily newspaper found that in 2017, e-scooters and roller skates combined caused 284 injuries and five deaths in France, a 23 percent increase on the previous year. The perception of e-scooters is that they are fast, silent and therefore dangerous, causing many accidents, and the need to legislate is urgent.
In an April 2023 referendum, voters in Paris chose to remove e-scooters from the city after the current vendor contracts expire. The ban applies to rental scooters which have been offered by several operators since 2018, although people will still be able to use privately owned contraptions.
Germany
In April 2019, low-power electric vehicles like e-scooters and segways were added to the regulatory list of vehicles allowed to circulate in the streets. The eKFV was enacted on 15 June 2019.
The regulation limits the maximum speed of these electric vehicles to 20 km/h and restricts them to cycle paths. Their operation requires a motor vehicle insurance and an insurance badge, but a driver's license is not required. Users have to be at least 14 years old.
Crash accident are under-reported (74% missing) when counted as declaration to police rather than to the hospital.
The same rules for operating an automobile while intoxicated also apply to electric kick scooters.
Ireland
The use of e-scooters and mono-wheels has exploded in Irish urban areas in recent years, with estimated more than 2,000 e-scooters regularly traveling the roads of Dublin.
Under existing road traffic legislation, the use of an e-scooter on public roads is not permitted. According to the Road Traffic Act 1961, all e-scooters are considered to be "mechanically propelled vehicles". Anyone using a mechanically propelled vehicle in a public place must have insurance, road tax, and a driving license. However, it is currently not possible to tax or insure e-scooters or electric skateboards.
In March 2019, e-scooter owners started reporting that the Irish police force, the Garda Síochána, had begun regularly seizing e-scooters on the grounds that the owner did not have insurance. This was despite a Freedom of Information request detailing that the Garda website displayed incorrect information to the public, detailing that e-scooters requiring human power to start would not be considered mechanically propelled vehicles and, as such, would fall outside the remit requiring insurance. The owner groups, such as eScoot.ie, have been publicly vocal, attracting media attention and urging e-scooter owners to sign a petition for lawmakers to legalize the public use of "electric rideables" in Ireland. Under growing pressure, the Minister for Transport Shane Ross asked the Road Safety Authority to research how e-scooters are regulated in other countries, particularly other EU member states. A decision is to be taken on whether or not to amend existing legislation.
In August 2019 the Road Safety Authority submitted a report on the use of e-scooters to Ross. The report is broadly in favour of e-scooters, however a number of significant safety concerns were raised. The Minister have announced a two-month public consultation starting on 1 September 2019. The main areas of the consultation cover what personal protective equipment should be used, what training should be provided, what safety or certification standards devices should meet, what age restrictions should apply and where the devices can be used publicly.
In February 2021 Communications Minister Eamon Ryan approved draft legislation which will "regularise" e-scooters and electric bikes as commonly accepted means of transport under proposed new vehicle category, to be known as "Powered Personal Transporters" (PPTs), which will not require road tax, insurance or driving license.
Japan
Japan is removing in July 2023 the requirement for escooter riders to have a driver's license. Scooters can be ridden on pavements where bicycles are allowed as long as they are slower than 6 kph and flash a green light.
Netherlands
Limited numbers of approved e-scooters by the Dutch Ministry of Infrastructure and Water Management are permitted to drive on public roads and categorised as 'Bijzondere Bromfiets'. The speed limit is 25km/h. Insurance, driving license and license plate not required for the approved models by the state.
New Zealand
E-scooters in New Zealand are classed as a 'Low-powered vehicle that does not require registration', provided that the output power is under 300 watts. They can therefore be ridden on footpaths, roads and separated cycleways. They cannot be ridden on paint-defined cycleways on the road. Helmets are not required, but recommended.
Norway
In Norway, e-scooters are classed as bicycles, and can therefore be ridden on footpaths, roads and separated cycleways as well as paint-defined cycleways on the road. Maximum speed is restricted to 20 km/h. Maximum weight of the e-scooter, including the battery, must not exceed 70 kg. Maximum width must not exceed 85 cm and maximum length is 120 cm. There is no age restriction or requirement to wear a helmet.
Helmets for children up to 15 years are mandatory since spring 2022.
Blood Alcohol Concentration (BAC) is limited to 0.2 gram per liter as for car drivers.
Poland
Following a court case, a new provision of the Road Traffic Act came into force as of 21 April 2019, whereby an e-scooter falls under the definition of a moped (power up to 4 kW, max speed 45 km/h). Therefore, such vehicles are not allowed to ride on the footpaths as well as bicycle lanes. However, due to the lack of homologation, it is not possible to register an e-scooter as a road vehicle, which makes it illegal for the use on the road. The legislators are now working on changes to the law to introduce the definition of the Personal Transport Device, which would allow e-scooters to be used on footpaths and bicycle lanes.
From May 20, 2021, the regulations on the traffic of e-scooters are in force. An e-scooter is an electric powered vehicle, two-axle, with a steering wheel, without a seat and without pedals, designed to be driven only by the rider on that vehicle.
To drive an e-scooter on the road by people aged 10 to 18, it is required to have the same qualifications as for cycling, i.e. a bicycle card or driving license of categories AM, A1, B1 or T. For people over 18 years, such a document is not required.
Singapore
E-scooters in Singapore are categorized as Personal Mobility Devices (PMD), and as such, are subjected to the Land Transport Authority's regulations. All e-scooter owners are required to register their devices with the Land Transport Authority and affix the registration number on their scooter. E-scooters that are not registered by 1 July 2019 will have their devices seized by the authorities and the offender would be liable for punishment.
E-scooters sold in Singapore have to comply with a strict set of regulations; maximum speed of , must not exceed 70 cm in width & must not weigh more than 20 kg. Retailers are allowed to sell non-compliant e-scooters however they have to indicate clearly that they can only be used on private property or for use overseas.
Unlike electric bicycles, e-scooters can only be ridden on footpaths and cycling paths. They are not allowed to be ridden on public roads.
Spain
E-scooters' recurring role in traffic accidents has led to a regulatory pushback in Spain. There have been reported 273 accidents, three of which were fatal in 2018. Spanish legislators are working on a regulation banning e-scooters from footpaths and limiting their speed to .
The first ever person hit by e-scooter died in Spain in August 2019. A 92-year-old woman fell and struck her head to the pavement when an e-scooter hit her, travelling at less than .
Spain is introducing technical standards and mandatory helmets.
Turkey
E-scooters can be used on cycle paths, and on urban roads without cycle paths where the speed limit is below 50 kph.
United Kingdom
Privately owned e-scooters are deemed to be Personal Light Electric Vehicles, subject to legal requirements regarding MOT testing, tax, and licensing. In practice they cannot be made to meet the requirements for road use, and they also may not be used on footways. In some trial areas from mid-2020 to November 2022, rental e-scooters may be ridden on roads and cycle lanes but not footways; riders must be 16 or over and have a driving licence. Using a phone, driving under the influence of alcohol, and other risks, are not allowed, as for other motor vehicles. Action is not usually taken against users of private scooters on roads and footways, but in December 2021 West Midlands Police announced that they had seized and destroyed 140 e-scooters. In July 2023, the police and crime commissioner for Kent called on police to seize and crush all e-scooters being ridden on public land.
In 2022 a woman riding a rental scooter erratically while over the legal limit for alcohol pleaded guilty to drink-driving. She had not known that it was an offence, but was fined, and banned from driving for 18 months.
Deaths
The first UK fatality involving an e-scooter occurred on 12 July 2019 when 35-year-old Emily Hartridge was killed in Battersea, London in a collision on a roundabout with a truck. London's cycling commissioner said that "new regulations must be put forward quickly" as e-scooters are "currently not safe—with no restrictions on speeds, no mandatory brakes and lights, and no rules on who can ride them and where".
The first death of a pedestrian hit by an e-scooter occurred on 8 June 2022, when the 71-year old victim died in hospital after being impacted by a 14-year old scooter-riding male on 2 June.
United States
Rules in the United States vary by state. Motorized scooters are often not street legal, as they cannot be tagged, titled, insured, and do not meet federal requirements for lights or mirrors. Particular localities may have further ordinances that limit the use of motorized scooters. The top speed of the average motorized scooter is around . Due to their small wheels, motorized scooters are not typically safe for street use as even the smallest bumps can cause an accident.
California, for example, requires that a person riding a motorized scooter on a street be 16 years of age or older, have a valid driver's license, be wearing a bicycle helmet, have no passengers, and otherwise follow the same rules of the road the same as cars do. The motorized scooter must have brakes, may not have handlebars raised above the operator's shoulders, and if ridden at night must have a headlight, a taillight, and side reflectors. A motorized scooter may not be operated on sidewalks or on streets if the posted speed limit is over unless in a Class II bicycle lane.
Michigan laws treat motorized scooters similarly to bicycles. They are typically allowed on sidewalks, bike lanes, and roads.
In Washington, D.C., motorized scooters are classified as Personal Mobility Devices, and are therefore not considered motor vehicles. This means there is no inspection, license, insurance, or registration required. Additionally, this means that motorized scooters are allowed on the sidewalks, and helmets are not required.
In Georgia, motorized scooters are considered Electric Personal Assistive Mobility Devices, meaning they can be used on sidewalks and highways where the speed limit is at most , or in the bike lane. The law also specifies that users of Electric Personal Assistive Mobility Devices, including motorized scooter riders, "have the same rights and duties as prescribed for pedestrians".
Scooter sharing companies have rules for operation printed on both the scooter and in the app, which includes instructions to not ride on the sidewalk. Given that the laws regarding motorized scooters vary from state to state, the scooter sharing instructions can differ from the local law.
Mechanics
Wheels and tires
Stand-up scooters may have solid tires, pneumatic tires with tubes, or tubeless pneumatic tires. Pneumatic tires offer benefits such as better shock absorption, adjustable tire pressure, and easier changes; however, they are prone to flats and require regular maintenance, making them ideal primarily for flat surfaces. Solid tires, often honeycomb in structure, have advantages such as a longer lifespan, puncture resistance, and low maintenance needs. However, they tend to have heavier weight and less shock absorption compared to air-filled tires. Sizes vary between and usually, and scooters with larger are available, for both road and off-road use. There are some with unusually wide tires especially for off-road use. Most of them use a steel or aluminum split rim.
Drive and transmissions
The simplest drive mechanism of stand-up scooters is the electric direct drive, where the motor directly drives the rear wheel. Some electric scooters have two motors, one for each wheel. Brushless motors can be extremely efficient this way, especially when regenerative braking is implemented. A large proportion of newer so-called "e-scooters" are designed this way.
When electric direct drive is not the rule, the simplest is the spindle drive, which puts an extension of the engine's output shaft, the spindle, in direct contact with the scooter's rear tire. To work correctly, the tire must have a clean, dry surface with which the spindle can effectively interact. Scooters with this type of direct transmission can be pull-started with the rear wheel off the ground, or "bump"-started by forcefully pushing them with the rear tire in contact with the ground.
Simple chain reduction drives are also used to transfer energy to the rear wheel, generally incorporating a type of centrifugal clutch to allow the engine to idle independently.
Belt reduction drives use the combination of wide flat "cog" belts and pulleys to transfer power to the rear wheel. Like chain drives, belt drives include a centrifugal clutch, but are more susceptible to breakage in off-road conditions.
Suspension
The suspension systems of stand-up scooters range from nothing at all, to simplistic spring based fork systems, to the complicated, dampened cam-link and C.I.D.L.I (Cantilevered Independent Dynamic Linkless Indespension) suspension mechanisms or a hybrid combination of wooden deck, coil spring, and dampers.
Brakes
Brake systems of kick scooters include disc brakes; magnetic brakes; and less efficient hydraulic brakes. Brakes can be placed on the front and/or back wheel(s). Many newer e-scooter models also have Kinetic Energy Regeneration System (KERS), which also acts as an electronic ABS system (E-ABS) on some models.
Gallery
Companies
| Technology | Motorized road transport | null |
6021465 | https://en.wikipedia.org/wiki/Biological%20pigment | Biological pigment | Biological pigments, also known simply as pigments or biochromes, are substances produced by living organisms that have a color resulting from selective color absorption. Biological pigments include plant pigments and flower pigments. Many biological structures, such as skin, eyes, feathers, fur and hair contain pigments such as melanin in specialized cells called chromatophores. In some species, pigments accrue over very long periods during an individual's lifespan.
Pigment color differs from structural color in that it is the same for all viewing angles, whereas structural color is the result of selective reflection or iridescence, usually because of multilayer structures. For example, butterfly wings typically contain structural color, although many butterflies have cells that contain pigment as well.
Biological pigments
See conjugated systems for electron bond chemistry that causes these molecules to have pigment.
Heme/porphyrin-based: chlorophyll, bilirubin, hemocyanin, hemoglobin, myoglobin
Light-emitting: luciferin
Carotenoids:
Hematochromes (algal pigments, mixes of carotenoids and their derivates)
Carotenes: alpha and beta carotene, lycopene, rhodopsin
Xanthophylls: canthaxanthin, zeaxanthin, lutein
Proteinaceous: phytochrome, phycobiliproteins
Psittacofulvins: a class of red and yellow pigments unique to parrots
Turacin and Turacoverdin: red and green pigments found in turacos and related species
Anthoxanthins: white color of some plants
Other: melanin, urochrome, flavonoids
Pigments in plants
The primary function of pigments in plants is photosynthesis, which uses the green pigment chlorophyll and several colorful pigments that absorb as much light energy as possible. Pigments are also known to play a role in pollination where pigment accumulation or loss can lead to floral color change, signaling to pollinators which flowers are rewarding and contain more pollen and nectar.
Plant pigments include many molecules, such as porphyrins, carotenoids, anthocyanins and betalains. All biological pigments selectively absorb certain wavelengths of light while reflecting others.
The principal pigments responsible are:
Chlorophyll is the primary pigment in plants; it is a chlorin that absorbs blue and red wavelengths of light while reflecting a majority of green. It is the presence and relative abundance of chlorophyll that gives plants their green color. All land plants and green algae possess two forms of this pigment: chlorophyll a and chlorophyll b. Kelps, diatoms, and other photosynthetic heterokonts contain chlorophyll c instead of b, while red algae possess only chlorophyll a. All chlorophylls serve as the primary means plants use to intercept light in order to fuel photosynthesis.
Carotenoids are red, orange, or yellow tetraterpenoids. During the process of photosynthesis, they have functions in light-harvesting (as accessory pigments), in photoprotection (energy dissipation via non-photochemical quenching as well as singlet oxygen scavenging for prevention of photooxidative damage), and also serve as protein structural elements. In higher plants, they also serve as precursors to the plant hormone abscisic acid.
Betalains are red or yellow pigments. Like anthocyanins they are water-soluble, but unlike anthocyanins they are synthesized from tyrosine. This class of pigments is found only in the Caryophyllales (including cactus and amaranth), and never co-occur in plants with anthocyanins. Betalains are responsible for the deep red color of beets.
Anthocyanins (literally "flower blue") are water-soluble flavonoid pigments that appear red to blue, according to pH. They occur in all tissues of higher plants, providing color in leaves, plant stem, roots, flowers, and fruits, though not always in sufficient quantities to be noticeable. Anthocyanins are most visible in the petals of flowers of many species.
Plants, in general, contain six ubiquitous carotenoids: neoxanthin, violaxanthin, antheraxanthin, zeaxanthin, lutein and β-carotene. Lutein is a yellow pigment found in fruits and vegetables and is the most abundant carotenoid in plants. Lycopene is the red pigment responsible for the color of tomatoes. Other less common carotenoids in plants include lutein epoxide (in many woody species), lactucaxanthin (found in lettuce), and alpha carotene (found in carrots).
A particularly noticeable manifestation of pigmentation in plants is seen with autumn leaf color, a phenomenon that affects the normally green leaves of many deciduous trees and shrubs whereby they take on, during a few weeks in the autumn season, various shades of red, yellow, purple, and brown.
Chlorophylls degrade into colorless tetrapyrroles known as nonfluorescent chlorophyll catabolites (NCCs).
As the predominant chlorophylls degrade, the hidden pigments of yellow xanthophylls and orange beta-carotene are revealed. These pigments are present throughout the year, but the red pigments, the anthocyanins, are synthesized de novo once roughly half of chlorophyll has been degraded. The amino acids released from degradation of light harvesting complexes are stored all winter in the tree's roots, branches, stems, and trunk until next spring when they are recycled to re‑leaf the tree.
Pigments in algae
Algae are very diverse photosynthetic organisms, which differ from plants in that they are aquatic organisms, they do not present vascular tissue and do not generate an embryo. However, both types of organisms share the possession of photosynthetic pigments, which absorb and release energy that is later used by the cell. These pigments in addition to chlorophylls, are phycobiliproteins, fucoxanthins, xanthophylls and carotenes, which serve to trap the energy of light and lead it to the primary pigment, which is responsible for initiating oxygenic photosynthesis reactions.
Algal phototrophs such as dinoflagellates use peridinin as a light harvesting pigment. While carotenoids can be found complexed within chlorophyll-binding proteins such as the photosynthetic reaction centers and light-harvesting complexes, they also are found within dedicated carotenoid proteins such as the orange carotenoid protein of cyanobacteria.
Pigments in bacteria
Bacteria produce pigments such as carotenoids, melanin, violacein, prodigiosin, pyocyanin, actinorhodin, and zeaxanthin. Cyanobacteria produce phycocyanin, phycoerythrin, scytonemin, chlorophyll a, chlorophyll d, and chlorophyll f. Purple sulfur bacteria produce bacteriochlorophyll a and bacteriochlorophyll b. In cyanobacteria, many other carotenoids exist such as canthaxanthin, myxoxanthophyll, synechoxanthin, and echinenone.
Pigments in animals
Pigmentation is used by many animals for protection, by means of camouflage, mimicry, or warning coloration. Some animals including fish, amphibians and cephalopods use pigmented chromatophores to provide camouflage that varies to match the background.
Pigmentation is used in signalling between animals, such as in courtship and reproductive behavior. For example, some cephalopods use their chromatophores to communicate.
The photopigment rhodopsin intercepts light as the first step in the perception of light.
Skin pigments such as melanin may protect tissues from sunburn by ultraviolet radiation.
However, some biological pigments in animals, such as heme groups that help to carry oxygen in the blood, are colored as a result of happenstance. Their color does not have a protective or signalling function.
Pea aphids (Acyrthosiphon pisum), two-spotted spider mites (Tetranychus urticae), and gall midges (family Cecidomyiidae) are the only known animals capable of synthesizing carotenoids. The presence of genes for synthesizing carotenoids in these arthropods has been attributed to independent horizontal gene transfer (HGT) events from fungi.
Diseases and conditions
A variety of diseases and abnormal conditions that involve pigmentation are in humans and animals, either from absence of or loss of pigmentation or pigment cells, or from the excess production of pigment.
Albinism is an inherited disorder characterized by total or partial loss of melanin. Humans and animals that suffer from albinism are called "albinistic" (the term "albino" is also sometimes used, but may be considered offensive when applied to people).
Lamellar ichthyosis, also called "fish scale disease", is an inherited condition in which one symptom is excess production of melanin. The skin is darker than normal, and is characterized by darkened, scaly, dry patches.
Melasma is a condition in which dark brown patches of pigment appear on the face, influenced by hormonal changes. When it occurs during a pregnancy, this condition is called the mask of pregnancy.
ocular pigmentation is an accumulation of pigment in the eye, and may be caused by latanoprost medication.
Vitiligo is a condition in which there is a loss of pigment-producing cells called melanocytes in patches of skin.
Pigments in marine animals
Carotenoids and carotenoproteins
Carotenoids are the most common group of pigments found in nature. Over 600 different kinds of carotenoids are found in animals, plants, and microorganisms.
Marine animals are incapable of making their own carotenoids and thus rely on plants for these pigments. Carotenoproteins are especially common among marine animals. These complexes are responsible for the various colors (red, purple, blue, green, etc.) to these marine invertebrates for mating rituals and camouflage. There are two main types of carotenoproteins: Type A and Type B. Type A has carotenoids (chromogen) which are stoichiometrically associated with a simple protein (glycoprotein). The second type, Type B, has carotenoids which are associated with a lipo protein and is usually less stable. While Type A is commonly found in the surface (shells and skins) of marine invertebrates, Type B is usually in eggs, ovaries, and blood. The colors and characteristic absorption of these carotenoprotein complexes are based upon the chemical binding of the chromogen and the protein subunits.
For example, the blue carotenoprotein, linckiacyanin has about 100-200 carotenoid molecules per every complex. In addition, the functions of these pigment-protein complexes also change their chemical structure as well. Carotenoproteins that are within the photosynthetic structure are more common, but complicated. Pigment-protein complexes that are outside of the photosynthetic system are less common, but have a simpler structure. For example, there are only two of these blue astaxanthin-proteins in the jellyfish, Velella velella, contains only about 100 carotenoids per complex.
A common carotenoid in animals is astaxanthin, which gives off a purple-blue and green pigment. Astaxanthin's color is formed by creating complexes with proteins in a certain order. For example, the crustochrin has approximately 20 astaxanthin molecules bonded with protein. When the complexes interact by exciton-exciton interaction, it lowers the absorbance maximum, changing the different color pigments.
In lobsters, there are various types of astaxanthin-protein complexes present. The first one is crustacyanin (max 632 nm), a slate-blue pigment found in the lobster's carapace. The second one is crustochrin (max 409), a yellow pigment which is found on the outer layer of the carapace. Lastly, the lipoglycoprotein and ovoverdin forms a bright green pigment that is usually present in the outer layers of the carapace and the lobster eggs.
Tetrapyrroles
Tetrapyrroles are the next most common group of pigments. They have four pyrrole rings, each ring consisting of C4H4NH. The main role of the tetrapyrroles is their connection in the biological oxidation process. Tetrapyrroles have a major role in electron transport and act as a replacement for many enzymes. They also have a role in the pigmentation of the marine organism's tissues.
Melanin
Melanin is a class of compounds that serves as a pigment with different structures responsible for dark, tan, yellowish / reddish pigments in marine animals. It is produced as the amino acid tyrosine is converted into melanin, which is found in the skin, hair, and eyes. Derived from aerobic oxidation of phenols, they are polymers.
There are several different types of melanins considering that they are an aggregate of smaller component molecules, such as nitrogen containing melanins. There are two classes of pigments: black and brown insoluble eumelanins, which are derived from aerobic oxidation of tyrosine in the presence of tyrosinase, and the alkali-soluble phaeomelanins which range from a yellow to red brown color, arising from the deviation of the eumelanin pathway through the intervention of cysteine and/or glutathione. Eumelanins are usually found in the skin and eyes. Several different melanins include melanoprotein (dark brown melanin that is stored in high concentrations in the ink sac of the cuttlefish Sepia Officianalis), echinoidea (found in sand dollars, and the hearts of sea urchins), holothuroidea (found in sea cucumbers), and ophiuroidea (found in brittle and snake stars). These melanins are possibly polymers which arise from the repeated coupling of simple bi-polyfunctional monomeric intermediates, or of high molecular weights. The compounds benzothiazole and tetrahydroisoquinoline ring systems act as UV-absorbing compounds.
Bioluminescence
The only light source in the deep sea, marine animals give off visible light energy called bioluminescence, a subset of chemiluminescence. This is the chemical reaction in which chemical energy is converted to light energy. It is estimated that 90% of deep-sea animals produce some sort of bioluminescence. Considering that a large proportion of the visible light spectrum is absorbed before reaching the deep sea, most of the emitted light from the sea-animals is blue and green. However, some species may emit a red and infrared light, and there has even been a genus that is found to emit yellow bioluminescence. The organ that is responsible for the emission of bioluminescence is known as photophores. This type is only present in squid and fish, and is used to illuminate their ventral surfaces, which disguise their silhouettes from predators. The uses of the photophores in the sea-animals differ, such as lenses for controlling intensity of color, and the intensity of the light produced. Squids have both photophores and chromatophores which controls both of these intensities. Another thing that is responsible for the emission of bioluminescence, which is evident in the bursts of light that jellyfish emit, start with a luciferin (a photogen) and ends with the light emitter (a photagogikon.) Luciferin, luciferase, salt, and oxygen react and combine to create a single unit called photo-proteins, which can produce light when reacted with another molecule such as Ca+. Jellyfish use this as a defense mechanism; when a smaller predator is attempting to devour a jellyfish, it will flash its lights, which would therefore lure a larger predator and chase the smaller predator away. It is also used as mating behavior.
In reef-building coral and sea anemones, they fluoresce; light is absorbed at one wavelength, and re-emitted at another. These pigments may act as natural sunscreens, aid in photosynthesis, serve as warning coloration, attract mates, warn rivals, or confuse predators.
Chromatophores
Chromatophores are color pigment changing cells that are directly stimulated by central motor neurons. They are primarily used for quick environmental adaptation for camouflaging. The process of changing the color pigment of their skin relies on a single highly developed chromatophore cell and many muscles, nerves, glial and sheath cells. Chromatophores contract and contain vesicles that stores three different liquid pigments. Each color is indicated by the three types of chromatophore cells: erythrophores, melanophores, and xanthophores. The first type is the erythrophores, which contains reddish pigments such as carotenoids and pteridines. The second type is the melanophores, which contains black and brown pigments such as the melanins. The third type is the xanthophores which contains yellow pigments in the forms of carotenoids. The various colors are made by the combination of the different layers of the chromatophores. These cells are usually located beneath the skin or scale the animals. There are two categories of colors generated by the cell – biochromes and schematochromes. Biochromes are colors chemically formed microscopic, natural pigments. Their chemical composition is created to take in some color of light and reflect the rest. In contrast, schematochromes (structural colors) are colors created by light reflections from a colorless surface and refractions by tissues. Schematochromes act like prisms, refracting and dispersing visible light to the surroundings, which will eventually reflect a specific combination of colors. These categories are determined by the movement of pigments within the chromatophores. The physiological color changes are short-term and fast, found in fishes, and are a result from an animal's response to a change in the environment. In contrast, the morphological color changes are long-term changes, occurs in different stages of the animal, and are due to the change of numbers of chromatophores. To change the color pigments, transparency, or opacity, the cells alter in form and size, and stretch or contract their outer covering.
Photo-protective pigments
Due to damage from UV-A and UV-B, marine animals have evolved to have compounds that absorb UV light and act as sunscreen. Mycosporine-like amino acids (MAAs) can absorb UV rays at 310-360 nm. Melanin is another well-known UV-protector. Carotenoids and photopigments both indirectly act as photo-protective pigments, as they quench oxygen free-radicals. They also supplement photosynthetic pigments that absorb light energy in the blue region.
Defensive role of pigments
It's known that animals use their color patterns to warn off predators, however it has been observed that a sponge pigment mimicked a chemical which involved the regulation of moulting of an amphipod that was known to prey on sponges. So whenever that amphipod eats the sponge, the chemical pigments prevents the moulting, and the amphipod eventually dies.
Environmental influence on color
Coloration in invertebrates varies based on the depth, water temperature, food source, currents, geographic location, light exposure, and sedimentation. For example, the amount of carotenoid a certain sea anemone decreases as we go deeper into the ocean. Thus, the marine life that resides on deeper waters is less brilliant than the organisms that live in well-lit areas due to the reduction of pigments. In the colonies of the colonial ascidian-cyanophyte symbiosis Trididemnum solidum, their colors are different depending on the light regime in which they live. The colonies that are exposed to full sunlight are heavily calcified, thicker, and are white. In contrast the colonies that live in shaded areas have more phycoerythrin (pigment that absorbs green) in comparison to phycocyanin (pigment that absorbs red), thinner, and are purple. The purple color in the shaded colonies are mainly due to the phycobilin pigment of the algae, meaning the variation of exposure in light changes the colors of these colonies.
Adaptive coloration
Aposematism is the warning coloration to signal potential predators to stay away. In many chromodorid nudibranchs, they take in distasteful and toxic chemicals emitted from sponges and store them in their repugnatorial glands (located around the mantle edge). Predators of nudibranchs have learned to avoid these certain nudibranchs based on their bright color patterns. Preys also protect themselves by their toxic compounds ranging from a variety of organic and inorganic compounds.
Physiological activities
Pigments of marine animals serve several different purposes, other than defensive roles. Some pigments are known to protect against UV (see photo-protective pigments.) In the nudibranch Nembrotha Kubaryana, tetrapyrrole pigment 13 has been found to be a potent antimicrobial agent. Also in this creature, tamjamines A, B, C, E, and F has shown antimicrobial, antitumor, and immunosuppressive activities.
Sesquiterpenoids are recognized for their blue and purple colors, but it has also been reported to exhibit various bioactivities such as antibacterial, immunoregulating, antimicrobial, and cytotoxic, as well as the inhibitory activity against cell division in the fertilized sea urchin and ascidian eggs. Several other pigments have been shown to be cytotoxic. In fact, two new carotenoids that were isolated from a sponge called Phakellia stelliderma showed mild cytotoxicity against mouse leukemia cells. Other pigments with medical involvements include scytonemin, topsentins, and debromohymenialdisine have several lead compounds in the field of inflammation, rheumatoid arthritis and osteoarthritis respectively. There's evidence that topsentins are potent mediators of immunogenic inflation, and topsentin and scytonemin are potent inhibitors of neurogenic inflammation.
Uses
Pigments may be extracted and used as dyes.
Pigments (such as astaxanthin and lycopene) are used as dietary supplements.
| Biology and health sciences | Biochemistry and molecular biology | null |
6025205 | https://en.wikipedia.org/wiki/Groundwater%20recharge | Groundwater recharge | Groundwater recharge or deep drainage or deep percolation is a hydrologic process, where water moves downward from surface water to groundwater. Recharge is the primary method through which water enters an aquifer. This process usually occurs in the vadose zone below plant roots and is often expressed as a flux to the water table surface. Groundwater recharge also encompasses water moving away from the water table farther into the saturated zone. Recharge occurs both naturally (through the water cycle) and through anthropogenic processes (i.e., "artificial groundwater recharge"), where rainwater and/or reclaimed water is routed to the subsurface.
The most common methods to estimate recharge rates are: chloride mass balance (CMB); soil physics methods; environmental and isotopic tracers; groundwater-level fluctuation methods; water balance (WB) methods (including groundwater models (GMs)); and the estimation of baseflow (BF) to rivers.
Processes
Diffused or focused mechanisms
Groundwater recharge can occur through diffuse or focused mechanisms. Diffuse recharge occurs when precipitation infiltrates through the soil to the water table, and is by definition distributed over large areas. Focused recharge occurs where water leaks from surface water sources (rivers, lakes, wadis, wetlands) or land surface depressions, and generally becomes more dominant with aridity.
Natural recharge
Water is recharged naturally by rain and snow melt and to a smaller extent by surface water (rivers and lakes). Recharge may be impeded somewhat by human activities including paving, development, or logging. These activities can result in loss of topsoil resulting in reduced water infiltration, enhanced surface runoff and reduction in recharge. Use of groundwater, especially for irrigation, may also lower the water tables. Groundwater recharge is an important process for sustainable groundwater management, since the volume-rate abstracted from an aquifer in the long term should be less than or equal to the volume-rate that is recharged.
Recharge can help move excess salts that accumulate in the root zone to deeper soil layers, or into the groundwater system. Tree roots increase water saturation into groundwater reducing water runoff. Flooding temporarily increases river bed permeability by moving clay soils downstream, and this increases aquifer recharge.
Wetlands
Wetlands help maintain the level of the water table and exert control on the hydraulic head. This provides force for groundwater recharge and discharge to other waters as well. The extent of groundwater recharge by a wetland is dependent upon soil, vegetation, site, perimeter to volume ratio, and water table gradient. Groundwater recharge occurs through mineral soils found primarily around the edges of wetlands. The soil under most wetlands is relatively impermeable. A high perimeter to volume ratio, such as in small wetlands, means that the surface area through which water can infiltrate into the groundwater is high. Groundwater recharge is typical in small wetlands such as prairie potholes, which can contribute significantly to recharge of regional groundwater resources. Researchers have discovered groundwater recharge of up to 20% of wetland volume per season.
Artificial groundwater recharge
Managed aquifer recharge (MAR) strategies to augment freshwater availability include streambed channel modification, bank filtration, water spreading and recharge wells. A facility in Orange County, California cleans and injects 100 million gallons per day; or 90 billion gallons per year.
Artificial groundwater recharge is becoming increasingly important in India, where over-pumping of groundwater by farmers has led to underground resources becoming depleted. In 2007, on the recommendations of the International Water Management Institute, the Indian government allocated to fund dug-well recharge projects (a dug-well is a wide, shallow well, often lined with concrete) in 100 districts within seven states where water stored in hard-rock aquifers had been over-exploited. Another environmental issue is the disposal of waste through the water flux such as dairy farms, industrial, and urban runoff.
Pollution in stormwater run-off collects in retention basins. Concentrating degradable contaminants can accelerate biodegradation. However, where and when water tables are high this affects appropriate design of detention ponds, retention ponds and rain gardens.
Depression-focused recharge
If water falls uniformly over a field such that field capacity of the soil is not exceeded, then negligible water percolates to groundwater. If instead water puddles in low-lying areas, the same water volume concentrated over a smaller area may exceed field capacity resulting in water that percolates down to recharge groundwater. The larger the relative contributing runoff area is, the more focused infiltration is. The recurring process of water that falls relatively uniformly over an area, flowing to groundwater selectively under surface depressions is depression focused recharge. Water tables rise under such depressions.
Depression focused groundwater recharge can be very important in arid regions. More rain events are capable of contributing to groundwater supply.
Depression focused groundwater recharge also profoundly effects contaminant transport into groundwater. This is of great concern in regions with karst geological formations because water can eventually dissolve tunnels all the way to aquifers, or otherwise disconnected streams. This extreme form of preferential flow, accelerates the transport of contaminants and the erosion of such tunnels. In this way depressions intended to trap runoff water—before it flows to vulnerable water resources—can connect underground over time. Cavitation of surfaces above into the tunnels, results in potholes or caves.
Deeper ponding exerts pressure that forces water into the ground faster. Faster flow dislodges contaminants otherwise adsorbed on soil and carries them along. This can carry pollution directly to the raised water table below and into the groundwater supply. Thus, the quality of water collecting in infiltration basins is of special concern.
Estimation methods
Rates of groundwater recharge are difficult to quantify. This is because other related processes, such as evaporation, transpiration (or evapotranspiration) and infiltration processes must first be measured or estimated to determine the balance. There are no widely applicable method available that can directly and accurately quantify the volume of rainwater that reaches the water table.
The most common methods to estimate recharge rates are: chloride mass balance (CMB); soil physics methods; environmental and isotopic tracers; groundwater-level fluctuation methods; water balance (WB) methods (including groundwater models (GMs)); and the estimation of baseflow (BF) to rivers.
Regional, continental and global estimates of recharge commonly derive from global hydrological models.
Physical
Physical methods use the principles of soil physics to estimate recharge. The direct physical methods are those that attempt to actually measure the volume of water passing below the root zone. Indirect physical methods rely on the measurement or estimation of soil physical parameters, which along with soil physical principles, can be used to estimate the potential or actual recharge. After months without rain the level of the rivers under humid climate is low and represents solely drained groundwater. Thus, the recharge can be calculated from this base flow if the catchment area is already known.
Chemical
Chemical methods use the presence of relatively inert water-soluble substances, such as an isotopic tracer or chloride, moving through the soil, as deep drainage occurs.
Numerical models
Recharge can be estimated using numerical methods, using such codes as Hydrologic Evaluation of Landfill Performance, UNSAT-H, SHAW (short form of Simultaneous Heat and Water Transfer model), WEAP, and MIKE SHE. The 1D-program HYDRUS1D is available online. The codes generally use climate and soil data to arrive at a recharge estimate and use the Richards equation in some form to model groundwater flow in the vadose zone.
Factors affecting groundwater recharge
Climate change
Urbanization
Further implications of groundwater recharge are a consequence of urbanization. Research shows that the recharge rate can be up to ten times higher in urban areas compared to rural regions. This is explained through the vast water supply and sewage networks supported in urban regions in which rural areas are not likely to obtain. Recharge in rural areas is heavily supported by precipitation, and this is the opposite for urban areas. Road networks and infrastructure within cities prevent surface water from percolating into the soil, resulting in most surface runoff entering storm drains for local water supply. As urban development continues to spread across various regions, groundwater recharge rates will increase relative to the existing rates of the previous rural region. A consequence of sudden influxes in groundwater recharge includes flash flooding. The ecosystem will have to adjust to the elevated groundwater surplus due to groundwater recharge rates. Additionally, road networks are less permeable compared to soil, resulting in higher amounts of surface runoff. Therefore, urbanization increases the rate of groundwater recharge and reduces infiltration, resulting in flash floods as the local ecosystem accommodates changes to the surrounding environment.
Adverse factors
Drainage
Impervious surfaces
Soil compaction
Groundwater pollution
| Physical sciences | Hydrology | Earth science |
1175965 | https://en.wikipedia.org/wiki/Phorusrhacidae | Phorusrhacidae | Phorusrhacids, colloquially known as terror birds, are an extinct family of large carnivorous, mostly flightless birds that were among the largest apex predators in South America during the Cenozoic era. Their definitive fossil records range from the Middle Eocene to the Late Pleistocene around , though some specimens suggest that they were present since the Early Eocene.
They ranged in height from . One of the largest specimens from the Early Pleistocene of Uruguay, possibly belonging to Devincenzia, would have weighed up to . Their closest modern-day relatives are believed to be the seriemas. Titanis walleri, one of the larger species, is known from Texas and Florida in North America. This makes the phorusrhacids the only known large South American predator to migrate north in the Great American Interchange that followed the formation of the Isthmus of Panama land bridge (the main pulse of the interchange began about 2.6 Ma ago; Titanis at 5 Ma was an early northward migrant).
It was once believed that T. walleri became extinct in North America around the time of the arrival of humans, but subsequent datings of Titanis fossils provided no evidence for their survival after 1.8 Ma. However, reports from Uruguay of new findings of phorusrachids such as a specimen of Psilopterus dating to 96,040 ± 6,300 years ago would imply that phorusrhacids survived in South America until the late Pleistocene.
Phorusrhacids may have even made their way into Africa and Europe, if the genus Lavocatavis from Algeria and Eleutherornis from France and Switzerland are included. However, the taxonomic placement of both taxa within phorusrhacids are considered highly questionable, and their remains are too fragmentary to be included in phylogenetic analyses. Possible specimens have also been discovered from the La Meseta Formation of Seymour Island, Antarctica, suggesting that this group had a wider geographical range in the Paleogene.
The closely related bathornithids occupied a similar ecological niche in North America across the Eocene to Early Miocene; some, like Paracrax, were similar in size to the largest phorusrhacids. At least one analysis recovers Bathornis as sister taxa to phorusrhacids, on the basis of shared features in the jaws and coracoid, though this has been seriously contested, as these might have evolved independently for the same carnivorous, flightless lifestyle.
Description
The neck can be divided into three main regions. In the higher regions of the neck, the phorusrhacid has bifurcate neural spines (BNS), while it has high neural spines in its lower regions. This suggests that the phorusrhacid had a highly flexible and developed neck allowing it to carry its heavy head and strike with terrifying speed and power. Although the phorusrhacid externally looks like it has a short neck, its flexible skeletal neck structure proves that it could expand farther beyond the expected reach and intimidate its prey using its height, allowing it to strike more easily. Once stretched out into its full length in preparation for a downward strike, its developed neck muscles and heavy head could produce enough momentum and power to cause fatal damage to the terror bird's prey.
Kelenken guillermoi, from the Langhian stage of the Miocene epoch, some 15 million years ago, discovered in the Collón Curá Formation in Patagonia in 2006, represents the largest bird skull yet found. The fossil has been described as being a , nearly intact skull. The beak is roughly long and curves in a hook shape that resembles an eagle's beak. Most species described as phorusrhacid birds were smaller, tall, but the new fossil belongs to a bird that probably stood about tall. Scientists theorize that the large terror birds were extremely nimble and quick runners, able to reach speeds of . Examination of phorusrhacid habitats also indicates that phorusrhacids may have presented intense competition to predatory metatherian sparassodonts such as borhyaenids and thylacosmilids, causing the mammalian predators to choose forested habitats to avoid the more successful and aggressive avian predators on the open plains.
The feet of the phorusrhacids had four toes, the first of which, known as the hallux, was reduced and did not touch the ground, while the others, corresponding to the second, third and fourth toes, were kept on the ground. Analysis of the resistance of the toes based on biomechanical models of curved beams, in particular of the second toe and its nail claw, indicate that it was modified into a "sickle claw" and was relatively uniform in various species and said claw would be relatively curved and large, which implies the need to keep it elevated to avoid wear or breakage due to contact with the ground, which would be achieved with a well-developed extensor tubercle and soft tissue pads on the fingers. The second toe, which was shorter and had fewer phalanges, also had more resistance and would make it easier to hold the claw off the ground and retain prey, a compromise with its predatory function and movement on the run, as occurs with modern seriemas, although to a lesser degree of specialization than dromaeosaurid dinosaurs. This is further supported by footprints from the Late Miocene of the Río Negro Formation, showcasing a trackway made by a mid-to-large sized terror bird with functionally didactyl footprints, the inner toe with the sickle claw raised mostly off the ground akin to their Mesozoic counterparts.
Skull structure
In the past, these birds were thought to have high beaks, round orbits, and vaulted braincases though there was never enough empirical evidence to support this. However, new fossils have been discovered in Comallo, Argentina. These skulls reveal that the terror bird has a triangular dorsal view, a rostrum that is hooked and more than half the length of the actual skull, and a more compact caudal portion. The external nares and antorbital fenestras (areas found in the nose) were found to be more square than triangular. These all contribute to a skull that is more rectangular in view rather than triangular. The structure of the fossils also suggest that these birds may have been swifter than originally thought.
A skull from a smaller subspecies of this bird was also found recently. With this fossil, it was found that the internal structure of the beak is hollow and reinforced with thin-walled trabeculae. There is also an absence of both zona flexoria palatina and zona flexoria arcus jugalis, which are key features that relate to the evolution of cranial akinesis. The discovery of this skull allows for the establishment of primary osteological homologies, which are useful in comparative anatomy, functional morphology, and phylogenetic studies.
Palaeobiology
Most phorusrhacids were very fast runners. All members possessed a large, sharp beak, a powerful neck and sharp talons. However, even with these attributes, the phorusrhacids are often assumed to have preyed on relatively small animals (about the size of a rabbit) that could be dispatched with a minimum of struggle. This is because with the phorusrhacids' beak proportions, the jaw could not generate a great deal of bite force with which to kill the prey. This is disputable as many big-game hunting predators such as Smilodon, great white sharks and Allosaurus have weaker bite forces and often laterally weak skulls as adaptations towards, not away from, killing large prey, relying instead on the presence of a cutting edge, a wide gape made possible by the reduction of jaw musculature, and the driving force of the body or neck. Since phorusrhacids share many of the same adaptations, such as a large, laterally flattened skull with a sharp-edged beak and powerful neck musculature, it is possible that they were specialized predators of relatively large prey.
The bones of the beak were tightly fused together, making the beak more resilient to force from the front to back direction, thus suggesting that it could cause a great amount of harm through pecking as opposed to side-to-side head movements like shaking prey. Generally speaking, it is thought that a terror bird would use its feet to injure prey by kicking it, and to hold the prey down and dispatch by pecking at it with its large beak. Larger prey may also have been attacked by pecking and kicking, or by using the beak as a blade to strike at or slash vital organs.
It has been recently shown that at least some phorusrhacids like Andalgalornis, while very fast runners in a straight line, were poor at tight turns at speed, which contradicts the idea of phorusrhacids being agile predators of small prey.
Diet
All phorusrhacids are thought to have been carnivorous. The strong downwards curve from the tip of this beak suggests that it ripped the flesh from the body of other animals; many extant bird species with this feature are carnivorous. CT scans performed on the skull of a phorusrhacid reveal that the species would not have been able to shake its prey side to side, but rather exert significant downward force. Florentino Ameghino claimed in a letter to Édouard Trouessart that he had specimens from Argentina of "petrified masses preserving skeletons of large rodents, Interatheriidae [small notoungulates] and even Proterotheriidae [deer-sized litopterns], with all their bones crushed and corroded, piled on with no apparent order and forming a nearly spherical mass with the skull in the center" that resembled giant owl pellets, suggesting that phorusrhacids may have swallowed their prey whole and regurgitated the indigestible parts similar to owls. However, Ameghino never formally described these specimens and they have not yet been relocated, making it difficult to determine if they are phorusrhacid pellets. Fossilized pellets from northwestern Argentina have also been suggested to pertain to small phorusrhacids like Procariama.
Classification
The etymology of the name Phorusrhacidae is based on the type genus Phorusrhacos. When first described by Florentino Ameghino in 1887, the etymology of Phorusrhacos was not given. Current thinking is that the name is derived from a combination of the Greek words "phoros", which means bearer or bearing, and "rhakos", which translates to wrinkles, scars or rents. Researchers have compared Phorusrhacidae with the living families of Cariamidae and Sagittariidae, but their differences in body mass are too drastic and, thus, one cannot overly depend on these living families for answers.
During the early Cenozoic, after the extinction of the non-bird dinosaurs, mammals underwent an evolutionary diversification, and some bird groups around the world developed a tendency towards gigantism; this included the Gastornithidae, the Dromornithidae, the Palaeognathae, and the Phorusrhacidae. Phorusrhacids are an extinct group within Cariamiformes, the only living members of which are the two species of seriemas in the family Cariamidae. While they are the most taxon-rich group within Cariamiformes, the interrelationships between phorusrhacids are unclear due to the incompleteness of their remains. A lineage of related predatory birds, the bathornithids, occupied North America prior to the arrival of phorusrhacids, living from the Eocene to Miocene and filled a similar niche to phorusrhacids. Only one genus belongs in the family, Bathornis, according to a 2016 analysis by paleontologist Gerald Mayr, who noted that Bathornis was more lightly built, with longer limbs proportionally and skulls more akin to those of Cariama.
Phylogenetic analysis of Cariamiformes and their relatives according to Mayr (2016) in his redescription of Bathornis: A 2024 study finds Bathornis as closer to seriemas than phorusrhacids were.
Following the revision by Alvarenga and Höfling (2003), there are now 5 subfamilies, containing 14 genera and 18 species: These species were the product of adaptive radiation. The following classification is based on LaBarge, Garderner & Organ (2024), and taxa identified as incertae sedis were all excluded from phylogenetic analysis in their study (except for Brontornis):
Family Phorusrhacidae
Incertae sedis
Genus ?Lavocatavis – Middle Eocene Glib Zegdou Formation of Algeria (likely more related to a possible paleognath Eremopezus)
Genus ?Patagorhacos – Early Miocene Chichinales Formation of Rio Negro Province, Argentina.
Genus ?Paleopsilopterus (Lower Eocene (Itaboraian) Itaboraí Formation of Itaboraí, Brazil) (identity as a phorusrhacid dubious)
Genus ?Brontornis (Early to Middle Miocene (Santacrucian-Laventan) Santa Cruz and Monte León Formations, Argentina) — gigantic species, standing on average high. Placement in Phorusrhacidae and/or monophyly disputed.
Genus ?Eleutherornis- Middle Eocene (Bartonian) of Rhône, France and Baselland, Switzerland (a cariamiform, probably more related to Strigogyps)
Subfamily Physornithinae — equivalent to Brontornithinae, if Brontornis is included within the family
Genus Paraphysornis (Late Oligocene to Early Miocene (Deseadan) Tremembé Formation of São Paulo State, Brazil)
Genus Physornis (Middle to Late Oligocene (Deseadan) Sarmiento Formation of Santa Cruz Province, Argentina)
Subfamily Phorusrhacinae — giant species high (Kelenken up to high), but somewhat slender and decidedly more nimble than the Brontornithinae
Genus Devincenzia (Miocene to Early Pliocene, possibly up to Early Pleistocene)
Genus Kelenken (Middle Miocene (Colloncuran) Collón Curá Formation of Río Negro Province, Argentina; largest known phorusrhacid)
Genus Phorusrhacos (Early to Middle Miocene (Santacrucian) Santa Cruz Formation of Argentina)
Genus Titanis (Early Pliocene to Early Pleistocene (Blancan) of Florida, California, and Texas)
Subfamily Patagornithinae — intermediate sized and very nimble species, standing around high
Genus Patagornis (Early to Middle Miocene (Santacrucian-Laventan) Santa Cruz Formation of Santa Cruz Province, Argentina) – includes Morenomerceraria, Palaeociconia, Tolmodus
Genus Andrewsornis (Middle to Late Oligocene (Deseadan) Agua de la Piedra Formation of southern Argentina)
Genus Andalgalornis (Late Miocene to Early Pliocene (Huayquerian) Ituzaingó Formation of northwestern Argentina)
Subfamily Psilopterinae — small species, standing high
Genus Psilopterus (Middle Oligocene (Deseadan) Santa Cruz Formation and Late Miocene (Chasicoan) Arroyo Chasicó Formation of southern and eastern Argentina respectively) (Possible Late Pleistocene (Lujanian) records from Uruguay)
Genus Procariama (Late Miocene to Early Pliocene (Huayquerian-Montehermosan) Cerro Azul and Andalhualá Formations of Catamarca Province, Argentina)
Subfamily Mesembriornithinae — medium-sized species, standing high
Genus Mesembriornis (Late Miocene to Late Pliocene (Montehermosan) Monte Hermoso Formation of Argentina)
Genus Llallawavis (Late Pliocene (Chapadmalalan) Playa Los Lobos Allo Formation of northeastern Argentina)
Alvarenga and Höfling did not include the Ameghinornithidae from Europe in the phorusrhacoids; these have meanwhile turned out to be more basal members of Cariamae. Though traditionally considered as members of the Gruiformes, based on both morphological and genetic studies (the latter being based on the seriema) Cariamiformes may belong to a separate group of birds, Australaves, and their closest living relatives, according to nuclear sequence studies, are a clade consisting of Falconidae, Psittaciformes and Passeriformes.
The following cladogram follows the analysis of Degrange and colleagues, 2015:
Extinction
During the Miocene and early Pliocene epochs, there was an increase in the phorusrhacid population size in South America, suggesting that, in that time frame, the various species flourished as predators in the savanna environment.
With the emergence of the Isthmus of Panama 2.7 million years ago, carnivorous dogs, bears, and cats from North America were able to cross into South America, increasing competition. (They had been preceded by procyonids as early as 7.3 million years ago.) The population of phorusrhacids declined thereafter according to older hypotheses, suggesting that competition with newly arrived predators was a major contributor to their extinction. Similar ideas have been considered for sparassodonts and for South America's terrestrial sebecid crocodilians.
However, the role of competitive displacement in South American predator lineages has been questioned by some researchers. The timing of turnover events and the decline of South American predators do not correlate well with the arrival of large carnivores like canids or sabretooths (although they do correlate well with the earlier-arriving procyonids, which evolved to large body size in South America, but these were omnivorous), with native South American predator lineages (including most phorusrhacids and all sparassodonts and sebecids) dying out well before the arrival of most larger placental carnivores. Bathornithids, which were similar in ecology and are likely close relatives of phorusrhacids, existed entirely within North America during part of the Cenozoic and competed successfully for a time with large carnivorans such as nimravids, before becoming extinct in the Early Miocene, about 20 million years ago. The phorusrhacid Titanis expanded northward into southern North America during the Interchange and coexisted for several million years with large canids and big cats like Xenosmilus, before its extinction about 1.8 million years ago.
There were some suggestions that phorusrhacids, like the majority of Pleistocene megafauna, were killed off by human activity such as hunting or habitat change. This idea is no longer considered valid, as improved dating on Titanis specimens show that the last phorusrhacids went extinct over one million years before humans arrived. However, several fossil finds of smaller forms have been described from the late Pleistocene of Uruguay in South America. Psilopterus may have been present until 96,040 ± 6,300 years ago (maximum age obtained from the bottom of the fossil-containing stratum), which would extend the existence of the smaller members of this group of avian predators considerably. Another unidentified smaller type which may be a possible psilopterine from the La Paz Local Fauna of Uruguay has also been dated to the late Pleistocene, perhaps 17,620 ± 100 years ago based on radiocarbon analysis using accelerator mass spectrometry (AMS) for the molar enamel samples of a proboscidean from the same site, but the validity of this previous radiocarbon dating has been considered highly questionable due to the enamel's lack of collagen; the tibia of Macrauchenia patachonica from the same site has been more precisely dated to a mean value of approximately 21,600 ± 1,000 years ago based on gamma spectrometry and radiocarbon dating.
| Biology and health sciences | Prehistoric birds | Animals |
1177019 | https://en.wikipedia.org/wiki/Copper%28II%29%20oxide | Copper(II) oxide | Copper(II) oxide or cupric oxide is an inorganic compound with the formula CuO. A black solid, it is one of the two stable oxides of copper, the other being Cu2O or copper(I) oxide (cuprous oxide). As a mineral, it is known as tenorite, or sometimes black copper. It is a product of copper mining and the precursor to many other copper-containing products and chemical compounds.
Production
It is produced on a large scale by pyrometallurgy, as one stage in extracting copper from its ores. The ores are treated with an aqueous mixture of ammonium carbonate, ammonia, and oxygen to ultimately give copper(II) ammine complex carbonates, such as . After extraction from the residues and after separation from iron, lead, etc. impurities, the carbonate salt is decomposed with steam to give CuO.
It can be formed by heating copper in air at around 300–800 °C:
For laboratory uses, copper(II) oxide is conveniently prepared by pyrolysis of copper(II) nitrate or basic copper(II) carbonate:
(180°C)
Dehydration of cupric hydroxide has also been demonstrated:
Reactions
Copper(II) oxide reacts with mineral acids such as hydrochloric acid, sulfuric acid, and nitric acid to give the corresponding hydrated copper(II) salts:
CuO + 2 HNO3 → Cu(NO3)2 + H2O
CuO + 2 HCl → CuCl2 + H2O
CuO + H2SO4 → CuSO4 + H2O
In presence of water it reacts with concentrated alkali to form the corresponding cuprate salts:
2 NaOH + CuO + H2O → Na2[Cu(OH)4]
It can also be reduced to copper metal using hydrogen, carbon monoxide, and carbon:
CuO + H2 → Cu + H2O
CuO + CO → Cu + CO2
2 CuO + C → 2Cu + CO2
When cupric oxide is substituted for iron oxide in thermite the resulting mixture is a low explosive, not an incendiary.
Structure and physical properties
Copper(II) oxide belongs to the monoclinic crystal system. The copper atom is coordinated by 4 oxygen atoms in an approximately square planar configuration.
The work function of bulk CuO is 5.3 eV.
Uses
As a significant product of copper mining, copper(II) oxide is the starting point for the production of many other copper salts. For example, many wood preservatives are produced from copper oxide.
Cupric oxide is used as a pigment in ceramics to produce blue, red, and green, and sometimes gray, pink, or black glazes.
It is incorrectly used as a dietary supplement in animal feed. Due to low bioactivity, negligible copper is absorbed.
It is used when welding with copper alloys.
A copper oxide electrode formed part of the early battery type known as the Edison–Lalande cell. Copper oxide was also used in a lithium battery type (IEC 60086 code "G").
Pyrotechnics and fireworks
Used as moderate blue coloring agent in blue flame compositions with additional chlorine donors and oxidizers such as chlorates and perchlorates. Providing oxygen it can be used as flash powder oxidizer with metal fuels such as magnesium, aluminium, or magnalium powder. Sometimes it is used in strobe effects and thermite compositions as crackling stars effect.
Similar compounds
An example of natural copper(I,II) oxide is the mineral paramelaconite, Cu+2Cu2+2O3.
| Physical sciences | Oxide salts | Chemistry |
1177234 | https://en.wikipedia.org/wiki/Potassium%20dichromate | Potassium dichromate | Potassium dichromate, , is a common inorganic chemical reagent, most commonly used as an oxidizing agent in various laboratory and industrial applications. As with all hexavalent chromium compounds, it is acutely and chronically harmful to health. It is a crystalline ionic solid with a very bright, red-orange color. The salt is popular in laboratories because it is not deliquescent, in contrast to the more industrially relevant salt sodium dichromate.
Chemistry
Production
Potassium dichromate is usually prepared by the reaction of potassium chloride on sodium dichromate. Alternatively, it can be also obtained from potassium chromate by roasting chromite ore with potassium hydroxide. It is soluble in water and in the dissolution process it ionizes:
Reaction
Potassium dichromate is an oxidising agent in organic chemistry, and is milder than potassium permanganate. It is used to oxidize alcohols. It converts primary alcohols into aldehydes and, under more forcing conditions, into carboxylic acids. In contrast, potassium permanganate tends to give carboxylic acids as the sole products. Secondary alcohols are converted into ketones. For example, menthone may be prepared by oxidation of menthol with acidified dichromate. Tertiary alcohols cannot be oxidized.
In an aqueous solution the color change exhibited can be used to test for distinguishing aldehydes from ketones. Aldehydes reduce dichromate from the +6 to the +3 oxidation state, changing color from orange to green. This color change arises because the aldehyde can be oxidized to the corresponding carboxylic acid. A ketone will show no such change because it cannot be oxidized further, and so the solution will remain orange.
When heated strongly, it decomposes with the evolution of oxygen.
When an alkali is added to an orange-red solution containing dichromate ions, a yellow solution is obtained due to the formation of chromate ions (). For example, potassium chromate is produced industrially using potash:
The reaction is reversible.
Treatment with cold sulfuric acid gives red crystals of chromic anhydride (chromium trioxide, CrO3):
On heating with concentrated acid, oxygen is evolved:
Uses
Potassium dichromate has few major applications, as the sodium salt is dominant industrially. The main use is as a precursor to potassium chrome alum, used in leather tanning.
Cleaning
Like other chromium(VI) compounds (chromium trioxide, sodium dichromate), potassium dichromate has been used to prepare "chromic acid" for cleaning glassware and etching materials. Because of safety concerns associated with hexavalent chromium, this practice has been largely discontinued.
Construction
It is used as an ingredient in cement in which it retards the setting of the mixture and improves its density and texture. This usage commonly causes contact dermatitis in construction workers.
Photography and printing
In 1839, Mungo Ponton discovered that paper treated with a solution of potassium dichromate was visibly tanned by exposure to sunlight, the discoloration remaining after the potassium dichromate had been rinsed out. In 1852, Henry Fox Talbot discovered that exposure to ultraviolet light in the presence of potassium dichromate hardened organic colloids such as gelatin and gum arabic, making them less soluble.
These discoveries soon led to the carbon print, gum bichromate, and other photographic printing processes based on differential hardening. Typically, after exposure, the unhardened portion was rinsed away with warm water, leaving a thin relief that either contained a pigment included during manufacture or was subsequently stained with a dye. Some processes depended on the hardening only, in combination with the differential absorption of certain dyes by the hardened or unhardened areas. Because some of these processes allowed the use of highly stable dyes and pigments, such as carbon black, prints with an extremely high degree of archival permanence and resistance to fading from prolonged exposure to light could be produced.
Dichromated colloids were also used as photoresists in various industrial applications, most widely in the creation of metal printing plates for use in photomechanical printing processes.
Chromium intensification or Photochromos uses potassium dichromate together with equal parts of concentrated hydrochloric acid diluted down to approximately 10% v/v to treat weak and thin negatives of black and white photograph roll. This solution reconverts the elemental silver particles in the film to silver chloride. After thorough washing and exposure to actinic light, the film can be redeveloped to its end-point yielding a stronger negative which is able to produce a more satisfactory print.
A potassium dichromate solution in sulfuric acid can be used to produce a reversal negative (that is, a positive transparency from a negative film). This is effected by developing a black and white film but allowing the development to proceed more or less to the end point. The development is then stopped by copious washing and the film then treated in the acid dichromate solution. This converts the silver metal to silver sulfate, a compound that is insensitive to light. After thorough washing and exposure to actinic light, the film is developed again allowing the previously unexposed silver halide to be reduced to silver metal. The results obtained can be unpredictable, but sometimes excellent results are obtained producing images that would otherwise be unobtainable. This process can be coupled with solarisation so that the end product resembles a negative and is suitable for printing in the normal way.
Cr(VI) compounds have the property of tanning animal proteins when exposed to strong light. This quality is used in photographic screen-printing.
In screen-printing a fine screen of bolting silk or similar material is stretched taut onto a frame similar to the way canvas is prepared before painting. A colloid sensitized with a dichromate is applied evenly to the taut screen. Once the dichromate mixture is dry, a full-size photographic positive is attached securely onto the surface of the screen, and the whole assembly exposed to strong light – times vary from 3 minutes to a half an hour in bright sunlight – hardening the exposed colloid. When the positive is removed, the unexposed mixture on the screen can be washed off with warm water, leaving the hardened mixture intact, acting as a precise mask of the desired pattern, which can then be printed with the usual screen-printing process.
Analytical reagent
Because it is non-hygroscopic, potassium dichromate is a common reagent in classical "wet tests" in analytical chemistry.
Ethanol determination
The concentration of ethanol in a sample can be determined by back titration with acidified potassium dichromate. Reacting the sample with an excess of potassium dichromate, all ethanol is oxidized to acetic acid:
Full reaction of converting ethanol to acetic acid:
The excess dichromate is determined by titration against sodium thiosulfate. Adding the amount of excess dichromate from the initial amount, gives the amount of ethanol present. Accuracy can be improved by calibrating the dichromate solution against a blank.
One major application for this reaction is in old police breathalyzer tests. When alcohol vapor makes contact with the orange dichromate-coated crystals, the color changes from Cr(VI) orange to Cr(III) green. The degree of the color change is directly related to the level of alcohol in the suspect's breath.
Silver test
When dissolved in an approximately 35% nitric acid solution it is called Schwerter's solution and is used to test for the presence of various metals, notably for determination of silver purity. Pure silver will turn the solution bright red, sterling silver will turn it dark red, low grade coin silver (0.800 fine) will turn brown (largely due to the presence of copper which turns the solution brown) and even green for 0.500 silver.
Brass turns dark brown, copper turns brown, lead and tin both turn yellow while gold and palladium do not change.
Sulfur dioxide test
Potassium dichromate paper can be used to test for sulfur dioxide, as it turns distinctively from orange to green. This is typical of all redox reactions where hexavalent chromium is reduced to trivalent chromium. Therefore, it is not a conclusive test for sulfur dioxide. The final product formed is Cr2(SO4)3.
Wood treatment
Potassium dichromate is used to stain certain types of wood by darkening the tannins in the wood. It produces deep, rich browns that cannot be achieved with modern color dyes. It is a particularly effective treatment on mahogany.
Natural occurrence
Potassium dichromate occurs naturally as the rare mineral lópezite. It has only been reported as vug fillings in the nitrate deposits of the Atacama Desert of Chile and in the Bushveld igneous complex of South Africa.
Safety
In 2005–06, potassium dichromate was the 11th-most-prevalent allergen in patch tests (4.8%).
Potassium dichromate is one of the most common causes of chromium dermatitis; chromium is highly likely to induce sensitization leading to dermatitis, especially of the hand and forearms, which is chronic and difficult to treat. Toxicological studies have further illustrated its highly toxic nature. With rabbits and rodents, concentrations as low as 14 mg/kg have shown a 50% fatality rate amongst test groups. Aquatic organisms are especially vulnerable if exposed, and hence responsible disposal according to the local environmental regulations is advised.
As with other Cr(VI) compounds, potassium dichromate is carcinogenic. The compound is also corrosive and exposure may produce severe eye damage or blindness. Human exposure further encompasses impaired fertility.
| Physical sciences | Metallic oxyanions | Chemistry |
1177467 | https://en.wikipedia.org/wiki/Telehealth | Telehealth | Telehealth is the distribution of health-related services and information via electronic information and telecommunication technologies. It allows long-distance patient and clinician contact, care, advice, reminders, education, intervention, monitoring, and remote admissions. Telemedicine is sometimes used as a synonym, or is used in a more limited sense to describe remote clinical services, such as diagnosis and monitoring. When rural settings, lack of transport, a lack of mobility, conditions due to outbreaks, epidemics or pandemics, decreased funding, or a lack of staff restrict access to care, telehealth may bridge the gap
as well as provide distance-learning; meetings, supervision, and presentations between practitioners; online information and health data management and healthcare system integration. Telehealth could include two clinicians discussing a case over video conference; a robotic surgery occurring through remote access; physical therapy done via digital monitoring instruments, live feed and application combinations; tests being forwarded between facilities for interpretation by a higher specialist; home monitoring through continuous sending of patient health data; client to practitioner online conference; or even videophone interpretation during a consult.
Telehealth versus telemedicine
Telehealth is sometimes discussed interchangeably with telemedicine, the latter being more common than the former. The Health Resources and Services Administration distinguishes telehealth from telemedicine in its scope, defining telemedicine only as describing remote clinical services, such as diagnosis and monitoring, while telehealth includes preventative, promotive, and curative care delivery. This includes the above-mentioned non-clinical applications, like administration and provider education.
The United States Department of Health and Human Services states that the term telehealth includes "non-clinical services, such as provider training, administrative meetings, and continuing medical education", and that the term telemedicine means "remote clinical services".
The World Health Organization uses telemedicine to describe all aspects of health care including preventive care. The American Telemedicine Association uses the terms telemedicine and telehealth interchangeably, although it acknowledges that telehealth is sometimes used more broadly for remote health not involving active clinical treatments.
eHealth is another related term, used particularly in the U.K. and Europe, as an umbrella term that includes telehealth, electronic medical records, and other components of health information technology.
Methods and modalities
Telehealth requires good Internet access by participants, usually in the form of a strong, reliable broadband connection, and broadband mobile communication technology of at least the fourth generation (4G) or long-term evolution (LTE) standard to overcome issues with video stability and bandwidth restrictions. As broadband infrastructure has improved, telehealth usage has become more widely feasible.
Healthcare providers often begin telehealth with a needs assessment which assesses hardships which can be improved by telehealth such as travel time, costs or time off work. Collaborators, such as technology companies can ease the transition.
Delivery can come within four distinct domains: live video (synchronous), store-and-forward (asynchronous), remote patient monitoring, and mobile health.
Store and forward
Store-and-forward telemedicine involves acquiring medical data (like medical images, biosignals etc.) and then transmitting this data to a doctor or medical specialist at a convenient time for assessment offline. It does not require the presence of both parties at the same time. Dermatology (cf: teledermatology), radiology, and pathology are common specialties that are conducive to asynchronous telemedicine. A properly structured medical record preferably in electronic form should be a component of this transfer. The 'store-and-forward' process requires the clinician to rely on a history report and audio/video information in lieu of a physical examination.
Remote monitoring
Remote monitoring, also known as self-monitoring or testing, enables medical professionals to monitor a patient remotely using various technological devices. This method is primarily used for managing chronic diseases or specific conditions, such as heart disease, diabetes mellitus, or asthma. These services can provide comparable health outcomes to traditional in-person patient encounters, supply greater satisfaction to patients, and may be cost-effective. Examples include home-based nocturnal dialysis and improved joint management.
Real-time interactive
Electronic consultations are possible through interactive telemedicine services which provide real-time interactions between patient and provider. Videoconferencing has been used in a wide range of clinical disciplines and settings for various purposes including management, diagnosis, counseling and monitoring of patients.
Videotelephony
Videotelephony comprises the technologies for the reception and transmission of audio-video signals by users at different locations, for communication between people in real-time.
At the dawn of the technology, videotelephony also included image phones which would exchange still images between units every few seconds over conventional POTS-type telephone lines, essentially the same as slow scan TV systems.
Currently, videotelephony is particularly useful to the deaf and speech-impaired who can use them with sign language and also with a video relay service, and well as to those with mobility issues or those who are located in distant places and are in need of telemedical or tele-educational services.
Categories
Emergency care
Common daily emergency telemedicine is performed by SAMU Regulator Physicians in France, Spain, Chile and Brazil. Aircraft and maritime emergencies are also handled by SAMU centres in Paris, Lisbon and Toulouse.
A recent study identified three major barriers to adoption of telemedicine in emergency and critical care units. They include:
Regulatory challenges: related to the difficulty and cost of obtaining licensure across multiple states, malpractice protection and privileges at multiple facilities
Financial barrier: lack of acceptance and reimbursement by government payers and some commercial insurance carriers, which places the investment burden squarely upon the hospital or healthcare system.
Cultural barriers: occurring from the lack of desire, or unwillingness, of some physicians to adapt clinical paradigms for telemedicine applications.
Emergency Telehealth is also gaining acceptance in the United States. There are several modalities currently being practiced that include but are not limited to TeleTriage, TeleMSE and ePPE.
An example of telehealth in the field is when EMS arrives on scene of an incident and is able to take an EKG that is then sent directly to a physician at the hospital to be read. Therefore, allowing instant care and management.
Telenursing
Telenursing refers to the use of telecommunications and information technology in order to provide nursing services in health care whenever a large physical distance exists between patient and nurse, or between any number of nurses. As a field it is part of telehealth, and has many points of contacts with other medical and non-medical applications, such as telediagnosis, teleconsultation, telemonitoring, etc.
Telenursing is achieving significant growth rates in many countries due to several factors: the preoccupation in reducing the costs of health care, an increase in the number of aging and chronically ill population, and the increase in coverage of health care to distant, rural, small or sparsely populated regions. Among its benefits, telenursing may help solve increasing shortages of nurses; to reduce distances and save travel time, and to keep patients out of hospital. A greater degree of job satisfaction has been registered among telenurses.
In Australia, during January 2014, Melbourne tech startup Small World Social collaborated with the Australian Breastfeeding Association to create the first hands-free breastfeeding Google Glass application for new mothers. The application, named Google Glass Breastfeeding app trial, allows mothers to nurse their baby while viewing instructions about common breastfeeding issues (latching on, posture etc.) or call a lactation consultant via a secure Google Hangout, who can view the issue through the mother's Google Glass camera. The trial was successfully concluded in Melbourne in April 2014, and 100% of participants were breastfeeding confidently.
Telepalliative care
Palliative care is an interdisciplinary medical caregiving approach aimed at optimizing quality of life and mitigating suffering among people with serious, complex, and often terminal illnesses. In the past, palliative care was a disease specific approach, but today the World Health Organization (WHO) takes a broader approach suggesting that palliative care should be applied as early as possible to any chronic and fatal illness. As in many aspects of healthcare, telehealth is increasingly being used in palliative care and is often referred to as telepalliative care. The types of technology applied in telepalliative care are typically telecommunication technologies, such as video conferencing or messaging for follow-up, or digital symptom assessments through digital questionnaires generating alerts to health care professionals. Telepalliative care has been shown to be a feasible approach to deliver palliative care among patients, caregivers and health care professionals. Telepalliative care can provide an added support system that enable patients to remain at home through self-reporting of symptoms and tailoring care to specific patients. Studies have shown that the use of telehealth in palliative care is mostly well received by patients, and that telepalliative care may improve access to health care professionals at home and enhance feelings of security and safety among patients receiving palliative care. Further, telepalliative care may enable more efficient utilization of healthcare resources, promotes collaboration between different levels of healthcare and makes healthcare professionals more responsive to changes in patients' condition.
Challenging aspects of the use of telehealth in palliative care have also been described. Generally, palliative care is a diverse medical specialty, involving interdisciplinary professionals from different professional traditions and cultures, delivering care to a heterogenous cohort of patients with diverse diseases, conditions and symptoms. This makes it a challenge to develop telehealth that is suitable for all patients and in all contexts of palliative care. Some of the barriers to telepalliative care relate to inflexible reporting of complex and fluctuating symptoms and circumstances using electronic questionnaires. Further, palliative care emphasizes a holistic approach that should address existential, spiritual and mental distress related to serious illness. However, few studies have included the self-reporting of existential or spiritual concerns, emotions, and well-being. Healthcare professionals may also be uncomfortable providing emotional or psychological care remotely. Palliative care has been characterized as high-touch rather than high-tech, limiting the interest in applying technological advancements when developing interventions. To optimize the advantages and minimize the challenges with the use of telehealth in home-based palliative care, future research should include users in the design and development process. Understanding the potential of telehealth to support therapeutic relationships between patients and health care professionals and being aware of the possible difficulties and tensions it may create are critical to its successful and acceptable use.
Telepharmacy
Telepharmacy is the delivery of pharmaceutical care via telecommunications to patients in locations where they may not have direct contact with a pharmacist. It is an instance of the wider phenomenon of telemedicine, as implemented in the field of pharmacy. Telepharmacy services include drug therapy monitoring, patient counseling, prior authorization and refill authorization for prescription drugs, and monitoring of formulary compliance with the aid of teleconferencing or videoconferencing. Remote dispensing of medications by automated packaging and labeling systems can also be thought of as an instance of telepharmacy. Telepharmacy services can be delivered at retail pharmacy sites or through hospitals, nursing homes, or other medical care facilities. This approach allows patients in remote or underserved areas to receive pharmacy services that would otherwise be unavailable to them, enhancing access to care and ensuring continuity in medication management. Health outcomes appear similar when pharmacy services are delivered by telepharmacy compared to traditional service delivery.
The term can also refer to the use of videoconferencing in pharmacy for other purposes, such as providing education, training, and management services to pharmacists and pharmacy staff remotely.
Telepsychiatry
Teledentistry
Teledentistry is the use of information technology and telecommunications for dental care, consultation, education, and public awareness in the same manner as telehealth and telemedicine.
Teleaudiology
Tele-audiology is the utilization of telehealth to provide audiological services and may include the full scope of audiological practice. This term was first used by Gregg Givens in 1999 in reference to a system being developed at East Carolina University in North Carolina, US.
Teleneurology
Teleneurology describes the use of mobile technology to provide neurological care remotely, including care for stroke, movement disorders like Parkinson's disease, seizure disorders (e.g., epilepsy), etc. The use of teleneurology gives us the opportunity to improve health care access for billions around the globe, from those living in urban locations to those in remote, rural locations. Evidence shows that individuals with Parkinson's disease prefer personal connection with a remote specialist to their local clinician. Such home care is convenient but requires access to and familiarity with internet. A 2017 randomized controlled trial of "virtual house calls" or video visits with individuals diagnosed with Parkinson disease evidences patient preference for the remote specialist vs their local clinician after one year. Teleneurology for patients with Parkison's disease is found to be cheaper than in person visits by reducing transportation and travel time A recent systematic review by Ray Dorsey et al. describes both the limitations and potential benefits of teleneurology to improve care for patients with chronic neurological conditions, especially in low-income countries. White, well educated and technologically savvy people are the biggest consumers of telehealth services for Parkinson's disease. as compared to ethnic minorities in the US.
Teleneurosurgery
Telemedicine in neurosurgery was historically primarily used for follow-up visits by patients that had to travel far to undergo surgery. In the last decade, telemedicine was also used for remote ICU rounding as well as prompt evaluation for acute ischemic stroke and administration of IV alteplase in conjunction with Neurology. From the onset of the COVID-19 pandemic, there was a rapid surge in the use of telemedicine across all divisions of neurosurgery: vascular, oncology, spine, and functional neurosurgery. Not only for follow-up visits, but it has gained popularity for seeing new patients or following established patients regardless of whether they underwent surgery. Telemedicine is not limited to direct patient care only; there are a number of new research groups and companies focused on using telemedicine for clinical trials involving patients with neurosurgical diagnoses.
Teleneuropsychology
Teleneuropsychology is the use of telehealth/videoconference technology for the remote administration of neuropsychological tests. Neuropsychological tests are used to evaluate the cognitive status of individuals with known or suspected brain disorders and provide a profile of cognitive strengths and weaknesses. Through a series of studies, there is growing support in the literature showing that remote videoconference-based administration of many standard neuropsychological tests results in test findings that are similar to traditional in-person evaluations, thereby establishing the basis for the reliability and validity of teleneuropsychological assessment.
Telenutrition
Telenutrition refers to the use of video conferencing/ telephony to provide online consultation by a nutritionist or dietician. Patient or clients upload their vital statistics, diet logs, food pictures etc. on a telenutrition portal which are then used by nutritionist or dietician to analyze their current health condition. Nutritionist or dietician can then set goals for their respective client/ patients and monitor their progress regularly by follow-up consultations.
Telenutrition portals can help people seek remote consultation for themselves and/or their family. This can be extremely helpful for elderly or bed ridden patients who can consult their dietician from comfort of their homes.
Telenutrition showed to be feasible and the majority of patients trusted the nutritional televisits, in place of the scheduled but not provided follow-up visits during the lockdown of the COVID-19 pandemic.
Telerehabilitation
Telerehabilitation (or e-rehabilitation) is the delivery of rehabilitation services over telecommunication networks and the Internet. Most types of services fall into two categories: clinical assessment (the patient's functional abilities in his or her environment), and clinical therapy. Some fields of rehabilitation practice that have explored telerehabilitation are: neuropsychology, speech–language pathology, audiology, occupational therapy, and physical therapy. Telerehabilitation can deliver therapy to people who cannot travel to a clinic because the patient has a disability or because of travel time. Telerehabilitation also allows experts in rehabilitation to engage in a clinical consultation at a distance.
Most telerehabilitation is highly visual. As of 2014, the most commonly used mediums are webcams, videoconferencing, phone lines, videophones and webpages containing rich web applications. The visual nature of telerehabilitation technology limits the types of rehabilitation services that can be provided. It is most widely used for neuropsychological rehabilitation; fitting of rehabilitation equipment such as wheelchairs, braces or artificial limbs; and in speech-language pathology. Rich web applications for neuropsychological rehabilitation (aka cognitive rehabilitation) of cognitive impairment (from many etiologies) were first introduced in 2001. This endeavor has expanded as a teletherapy application for cognitive skills enhancement programs for school children. Tele-audiology (hearing assessments) is a growing application. Physical therapy and psychology interventions delivered via telehealth may result in similar outcomes as those delivered in-person for a range of health conditions.
Two important areas of telerehabilitation research are (1) demonstrating equivalence of assessment and therapy to in-person assessment and therapy, and (2) building new data collection systems to digitize information that a therapist can use in practice. Ground-breaking research in telehaptics (the sense of touch) and virtual reality may broaden the scope of telerehabilitation practice, in the future.
In the United States, the National Institute on Disability and Rehabilitation Research's (NIDRR) supports research and the development of telerehabilitation. NIDRR's grantees include the "Rehabilitation Engineering and Research Center" (RERC) at the University of Pittsburgh, the Rehabilitation Institute of Chicago, the State University of New York at Buffalo, and the National Rehabilitation Hospital in Washington DC. Other federal funders of research are the Veterans Health Administration, the Health Services Research Administration in the US Department of Health and Human Services, and the Department of Defense. Outside the United States, excellent research is conducted in Australia and Europe.
Only a few health insurers in the United States, and about half of Medicaid programs, reimburse for telerehabilitation services. If the research shows that teleassessments and teletherapy are equivalent to clinical encounters, it is more likely that insurers and Medicare will cover telerehabilitation services.
In India, the Indian Association of Chartered Physiotherapists (IACP) provides telerehabilitation facilities. With the support and collaboration of local clinics and private practitioners and the Members IACP, IACP runs the facility, named Telemedicine. IACP has maintained an internet-based list of their members on their website, through which patients can make online appointments.
Teletrauma care
Telemedicine can be utilized to improve the efficiency and effectiveness of the delivery of care in a trauma environment. Examples include:
Telemedicine for trauma triage: using telemedicine, trauma specialists can interact with personnel on the scene of a mass casualty or disaster situation, via the internet using mobile devices, to determine the severity of injuries. They can provide clinical assessments and determine whether those injured must be evacuated for necessary care. Remote trauma specialists can provide the same quality of clinical assessment and plan of care as a trauma specialist located physically with the patient.
Telemedicine for intensive care unit (ICU) rounds: Telemedicine is also being used in some trauma ICUs to reduce the spread of infections. Rounds are usually conducted at hospitals across the country by a team of approximately ten or more people to include attending physicians, fellows, residents and other clinicians. This group usually moves from bed to bed in a unit discussing each patient. This aids in the transition of care for patients from the night shift to the morning shift, but also serves as an educational experience for new residents to the team. A new approach features the team conducting rounds from a conference room using a video-conferencing system. The trauma attending, residents, fellows, nurses, nurse practitioners, and pharmacists are able to watch a live video stream from the patient's bedside. They can see the vital signs on the monitor, view the settings on the respiratory ventilator, and/or view the patient's wounds. Video-conferencing allows the remote viewers two-way communication with clinicians at the bedside.
Telemedicine for trauma education: some trauma centers are delivering trauma education lectures to hospitals and health care providers worldwide using video conferencing technology. Each lecture provides fundamental principles, firsthand knowledge and evidenced-based methods for critical analysis of established clinical practice standards, and comparisons to newer advanced alternatives. The various sites collaborate and share their perspective based on location, available staff, and available resources.
Telemedicine in the trauma operating room: trauma surgeons are able to observe and consult on cases from a remote location using video conferencing. This capability allows the attending to view the residents in real time. The remote surgeon has the capability to control the camera (pan, tilt and zoom) to get the best angle of the procedure while at the same time providing expertise in order to provide the best possible care to the patient.
Telecardiology
ECGs, or electrocardiographs, can be transmitted using telephone and wireless. Willem Einthoven, the inventor of the ECG, actually did tests with transmission of ECG via telephone lines. This was because the hospital did not allow him to move patients outside the hospital to his laboratory for testing of his new device. In 1906 Einthoven came up with a way to transmit the data from the hospital directly to his lab.
Transmission of ECGs
One of the oldest known telecardiology systems for teletransmissions of ECGs was established in Gwalior, India in 1975 at GR Medical college by Ajai Shanker, S. Makhija, P.K. Mantri using an indigenous technique for the first time in India.
This system enabled wireless transmission of ECG from the moving ICU van or the patients home to the central station in ICU of the department of Medicine. Transmission using wireless was done using frequency modulation which eliminated noise. Transmission was also done through telephone lines. The ECG output was connected to the telephone input using a modulator which converted ECG into high frequency sound. At the other end a demodulator reconverted the sound into ECG with a good gain accuracy. The ECG was converted to sound waves with a frequency varying from 500 Hz to 2500 Hz with 1500 Hz at baseline.
This system was also used to monitor patients with pacemakers in remote areas. The central control unit at the ICU was able to correctly interpret arrhythmia. This technique helped medical aid reach in remote areas.
In addition, electronic stethoscopes can be used as recording devices, which is helpful for purposes of telecardiology. There are many examples of successful telecardiology services worldwide.
In Pakistan three pilot projects in telemedicine were initiated by the Ministry of IT & Telecom, Government of Pakistan (MoIT) through the Electronic Government Directorate in collaboration with Oratier Technologies (a pioneer company within Pakistan dealing with healthcare and HMIS) and PakDataCom (a bandwidth provider). Three hub stations through were linked via the Pak Sat-I communications satellite, and four districts were linked with another hub. A 312 Kb link was also established with remote sites and 1 Mbit/s bandwidth was provided at each hub. Three hubs were established: the Mayo Hospital (the largest hospital in Asia), JPMC Karachi and Holy Family Rawalpindi. These 12 remote sites were connected and on average of 1,500 patients were treated per month per hub. The project was still running smoothly after two years.
Wireless ambulatory ECG technology, moving beyond previous ambulatory ECG technology such as the Holter monitor, now includes smartphones and Apple Watches which can perform at-home cardiac monitoring and send the data to a physician via the internet.
Teleradiology
Teleradiology is the ability to send radiographic images (x-rays, CT, MR, PET/CT, SPECT/CT, MG, US...) from one location to another. For this process to be implemented, three essential components are required, an image sending station, a transmission network, and a receiving-image review station. The most typical implementation are two computers connected via the Internet. The computer at the receiving end will need to have a high-quality display screen that has been tested and cleared for clinical purposes. Sometimes the receiving computer will have a printer so that images can be printed for convenience.
The teleradiology process begins at the image sending station. The radiographic image and a modem or other connection are required for this first step. The image is scanned and then sent via the network connection to the receiving computer.
Today's high-speed broadband based Internet enables the use of new technologies for teleradiology: the image reviewer can now have access to distant servers in order to view an exam. Therefore, they do not need particular workstations to view the images; a standard personal computer (PC) and digital subscriber line (DSL) connection is enough to reach Keosys' central server. No particular software is necessary on the PC and the images can be reached from anywhere in the world.
Teleradiology is the most popular use for telemedicine and accounts for at least 50% of all telemedicine usage.
Telepathology
Telepathology is the practice of pathology at a distance. It uses telecommunications technology to facilitate the transfer of image-rich pathology data between distant locations for the purposes of diagnosis, education, and research. Performance of telepathology requires that a pathologist selects the video images for analysis and the rendering diagnoses. The use of "television microscopy", the forerunner of telepathology, did not require that a pathologist have physical or virtual "hands-on" involvement is the selection of microscopic fields-of-view for analysis and diagnosis.
A pathologist, Ronald S. Weinstein, M.D., coined the term "telepathology" in 1986. In an editorial in a medical journal, Weinstein outlined the actions that would be needed to create remote pathology diagnostic services. He, and his collaborators, published the first scientific paper on robotic telepathology. Weinstein was also granted the first U.S. patents for robotic telepathology systems and telepathology diagnostic networks. Weinstein is known to many as the "father of telepathology". In Norway, Eide and Nordrum implemented the first sustainable clinical telepathology service in 1989. This is still in operation, decades later. A number of clinical telepathology services have benefited many thousands of patients in North America, Europe, and Asia.
Telepathology has been successfully used for many applications including the rendering histopathology tissue diagnoses, at a distance, for education, and for research. Although digital pathology imaging, including virtual microscopy, is the mode of choice for telepathology services in developed countries, analog telepathology imaging is still used for patient services in some developing countries.
Teledermatology
Teledermatology allows dermatology consultations over a distance using audio, visual and data communication, and has been found to improve efficiency, access to specialty care, and patient satisfaction. Applications comprise health care management such as diagnoses, consultation and treatment as well as (continuing medical) education. The dermatologists Perednia and Brown were the first to coin the term teledermatology in 1995, where they described the value of a teledermatologic service in a rural area underserved by dermatologists.
Teleophthalmology
Teleophthalmology is a branch of telemedicine that delivers eye care through digital medical equipment and telecommunications technology. Today, applications of teleophthalmology encompass access to eye specialists for patients in remote areas, ophthalmic disease screening, diagnosis and monitoring; as well as distant learning. Teleophthalmology may help reduce disparities by providing remote, low-cost screening tests such as diabetic retinopathy screening to low-income and uninsured patients. In Mizoram, India, a hilly area with poor roads, between 2011 and 2015, teleophthalmology provided care to over 10,000 patients. These patients were examined by ophthalmic assistants locally but surgery was done on appointment after the patient images were viewed online by eye surgeons in the hospital 6–12 hours away. Instead of an average five trips for say, a cataract procedure, only one was required for surgery alone as even post-op care like removal of stitches and appointments for glasses was done locally. There were large cost savings in travel as well.
In the United States, some companies allow patients to complete an online visual exam and within 24 hours receive a prescription from an optometrist valid for eyeglasses, contact lenses, or both. Some US states such as Indiana have attempted to ban these companies from doing business.
Telesurgery
Remote surgery (also known as telesurgery) is the ability for a doctor to perform surgery on a patient even though they are not physically in the same location. It is a form of telepresence. Remote surgery combines elements of robotics, cutting-edge telecommunications such as high-speed data connections, telehaptics and elements of management information systems. While the field of robotic surgery is fairly well established, most of these robots are controlled by surgeons at the location of the surgery.
Remote surgery is remote work for surgeons, where the physical distance between the surgeon and the patient is immaterial. It promises to allow the expertise of specialized surgeons to be available to patients worldwide, without the need for patients to travel beyond their local hospital.
Remote surgery or telesurgery is performance of surgical procedures where the surgeon is not physically in the same location as the patient, using a robotic teleoperator system controlled by the surgeon. The remote operator may give tactile feedback to the user. Remote surgery combines elements of robotics and high-speed data connections. A critical limiting factor is the speed, latency and reliability of the communication system between the surgeon and the patient, though trans-Atlantic surgeries have been demonstrated.
Teleabortion
Telemedicine has been used globally to increase access to abortion care, specifically medical abortion, in environments where few abortion care providers exist or abortion is legally restricted. Clinicians are able to virtually provide counseling, review screening tests, observe the administration of an abortion medication, and directly mail abortion pills to people. In 2004, Women on Web (WoW), Amsterdam, started offering online consultations, mostly to people living in areas where abortion was legally restricted, informing them how to safely use medical abortion drugs to end a pregnancy. People contact the Women on Web service online; physicians review any necessary lab results or ultrasounds, mail mifepristone and misoprostol pills to people, then follow up through online communication. In the United States, medical abortion was introduced as a telehealth service in Iowa by Planned Parenthood of the Heartland in 2008 to allow a patient at one health facility to communicate via secure video with a health provider at another facility. In this model a person seeking abortion care must come to a health facility. An abortion care provider communicates with the person located at another site using clinic-to-clinic videoconferencing to provide medical abortion after screening tests and consultation with clinic staff. In 2018, the website Aid Access was launched by the founder of Women on Web, Rebecca Gomperts. It offers a similar service as Women on Web in the United States, but the medications are prescribed to an Indian pharmacy, then mailed to the United States.
The TelAbortion study conducted by Gynuity Health Projects, with special approval from the U.S. Food and Drug Administration (FDA), aims to increase access to medical abortion care without requiring an in-person visit to a clinic. This models was expanded during the COVID-19 pandemic and as of March 2020 exists in 13 U.S. states and has enrolled over 730 people in the study. The person receives counseling and instruction from an abortion care provider via videoconference from a location of their choice. The medications necessary for the abortion, mifepristone and misoprostol, are mailed directly to the person and they have a follow-up video consultation in 7–14 days. A systematic review of telemedicine abortion has found the practice to be safe, effective, efficient, and satisfactory.
In the United States, eighteen states require the clinician to be physically present during the administration of medications for abortion which effectively bans telehealth of medication abortion: five states explicitly ban telemedicine for medication abortion, while thirteen states require the prescriber (usually required to be a physician) to be physically present with the patient. In the UK, the Royal College of Obstetricians and Gynecologists approved a no-test protocol for medication abortion, with mifepristone available through a minimal-contact pick-up or by mail.
Other specialist care delivery
Telemedicine can facilitate specialty care delivered by primary care physicians according to a controlled study of the treatment of hepatitis C. Various specialties are contributing to telemedicine, in varying degrees.Other specialist conditions for which telemedicine has been used include perinatal mental health.
In light of the COVID-19 pandemic, primary care physicians have relied on telehealth to continue to provide care in outpatient settings. The transition to virtual health has been beneficial in providing patients access to care (especially care that does not require a physical exam e.g. medication changes, minor health updates) and avoid putting patients at risk of COVID-19. This included providing services to pediatric patients during the pandemic, where issues of last minute cancelation and rescheduling were frequently related to a lack of technicality and engagement, two factors often understudied in the literature.
Telemedicine has also been beneficial in facilitating medical education to students while still allowing for adequate social distancing during the COVID-19 pandemic. Many medical schools have shifted to alternate forms of virtual curriculum and are still able to engage in meaningful telehealth encounters with patients.
Medication assisted treatment (MAT) is the treatment of opioid use disorder (OUD) with medications, often in combination with behavioral therapy As a response to the COVID-19 pandemic the use of telemedicine has been granted by the Drug Enforcement Administration to start or maintain people OUD on buprenorphine (trade name Suboxone) via telemedicine without the need for an initial in-person examination. On March 31, 2020, QuickMD became the first national TeleMAT service in the United States to provide Medication-assisted Treatment with Suboxone online – without the need of an in-person visit; with others announcing to follow soon.
Major developments
In policy
Telehealth is a modern form of health care delivery. Telehealth breaks away from traditional health care delivery by using modern telecommunication systems including wireless communication methods. Traditional health is legislated through policy to ensure the safety of medical practitioners and patients. Consequently, since telehealth is a new form of health care delivery that is now gathering momentum in the health sector, many organizations have started to legislate the use of telehealth into policy. In New Zealand, the Medical Council has a statement about telehealth on their website. This illustrates that the medical council has foreseen the importance that telehealth will have on the health system and have started to introduce telehealth legislation to practitioners along with government.
Transition to mainstream
Traditional use of telehealth services has been for specialist treatment. However, there has been a paradigm shift and telehealth is no longer considered a specialist service. This development has ensured that many access barriers are eliminated, as medical professionals and patients are able to use wireless communication technologies to deliver health care. This is evident in rural communities. Rural residents typically have to travel to longer distances to access healthcare than urban counterparts due to physician shortages and healthcare facility closures in these areas. Telehealth eliminates this barrier as health professionals are able to conduct medical consultations through the use of wireless communication technologies. However, this process is dependent on both parties having internet access and comfort level with technology, which poses barriers for many low-income and rural communities.
Telehealth allows the patient to be monitored between physician office visits which can improve patient health. Telehealth also allows patients to access expertise which is not available in their local area. This remote patient monitoring ability enables patients to stay at home longer and helps avoid unnecessary hospital time. In the long-term, this could potentially result in less burdening of the healthcare system and consumption of resources.
During the COVID-19 pandemic, there were large increases in the use of telemedicine for primary care visits within the United States, increasing from an average of 1.4 million visits in Q2 of 2018 and 2019 to 35 million visits in Q2 2020, according to data from IQVIA. The telehealth market is expected to grow at 40% a year in 2021. Use of telemedicine by General Practitioners in the UK rose from 20 to 30% pre-COVID to almost 80% by the beginning of 2021. More than 70% of practitioners and patients were satisfied with this. Boris Johnson was said to have "piled pressure on GPs to offer more in-person consultations" supporting a campaign largely orchestrated by the Daily Mail. The Royal College of General Practitioners said that a patient "right" to have face-to-face appointments if they wished was "undeliverable".
Technology advancement
The technological advancement of wireless communication devices is a major development in telehealth. This allows patients to self-monitor their health conditions and to not rely as much on health care professionals. Furthermore, patients are more willing to stay on their treatment plans as they are more invested and included in the process as the decision-making is shared. Technological advancement also means that health care professionals are able to use better technologies to treat patients for example in maternal care and surgery. A 2023 study published in the Journal of the American College of Surgeons showed telemedicine as making a positive impact, with expectations exceeded for those physicians and patients who had consulted online for surgeries. Technological developments in telehealth are essential to improve health care, especially the delivery of healthcare services, as resources are finite along with an ageing population that is living longer.
Licensing
U.S. licensing and regulatory issues
Restrictive licensure laws in the United States require a practitioner to obtain a full license to deliver telemedicine care across state lines. Typically, states with restrictive licensure laws also have several exceptions (varying from state to state) that may release an out-of-state practitioner from the additional burden of obtaining such a license. A number of states require practitioners who seek compensation to frequently deliver interstate care to acquire a full license.
If a practitioner serves several states, obtaining this license in each state could be an expensive and time-consuming proposition. Even if the practitioner never practices medicine face-to-face with a patient in another state, he/she still must meet a variety of other individual state requirements, including paying substantial licensure fees, passing additional oral and written examinations, and traveling for interviews.
In 2008, the U.S. passed the Ryan Haight Act which required face-to-face or valid telemedicine consultations prior to receiving a prescription.
State medical licensing boards have sometimes opposed telemedicine; for example, in 2012 electronic consultations were illegal in Idaho, and an Idaho-licensed general practitioner was punished by the board for prescribing an antibiotic, triggering reviews of her licensure and board certifications across the country. Subsequently, in 2015 the state legislature legalized electronic consultations.
In 2015, Teladoc filed suit against the Texas Medical Board over a rule that required in-person consultations initially; the judge refused to dismiss the case, noting that antitrust laws apply to state medical boards.
Major implications and impacts
Telehealth allows multiple, varying disciplines to merge and deliver a potentially more uniform level of care, using technology. As telehealth proliferates mainstream healthcare, it challenges notions of traditional healthcare delivery. Some populations experience better quality, access and more personalized health care.
Health promotion
Telehealth can also increase health promotion efforts. These efforts can now be more personalised to the target population and professionals can extend their help into homes or private and safe environments in which patients of individuals can practice, ask and gain health information. Health promotion using telehealth has become increasingly popular in underdeveloped countries where there are very poor physical resources available. There has been a particular push toward mHealth applications as many areas, even underdeveloped ones have mobile phone and smartphone coverage.
In a 2015 article reviewing research on the use of a mobile health application in the United Kingdom, authors describe how a home-based application helped patients manage and monitor their health and symptoms independently. The mobile health application allows people to rapidly self-report their symptoms – 95% of patients were able to report their daily symptoms in less than 100 seconds, which is less than the 5 minutes (plus commuting) taken to measure vital signs by nurses in hospitals. Online applications allow patients to remain at home to keep track of the progression of their chronic illnesses. The downside of using mHealth applications is that not everyone, especially in developing countries, has daily access to internet or electronic devices.
In developed countries, health promotion efforts using telehealth have been met with some success. The Australian hands-free breastfeeding Google Glass application reported promising results in 2014. This application made in collaboration with the Australian Breastfeeding Association and a tech startup called Small World Social, helped new mothers learn how to breastfeed. Breastfeeding is beneficial to infant health and maternal health and is recommended by the World Health Organisation and health organisations all over the world. Widespread breastfeeding can prevent 820,000 infant deaths globally but the practice is often stopped prematurely or intents to do are disrupted due to lack of social support, know-how or other factors. This application gave mother's hands-free information on breastfeeding, instructions on how to breastfeed and also had an option to call a lactation consultant over Google Hangout. When the trial ended, all participants were reported to be confident in breastfeeding.
Health care quality and barriers to adoption
A scientific review indicates that, in general, outcomes of telemedicine are or can be as good as in-person care with health care use staying similar.
Advantages of the nonexclusive adoption of already existing telemedicine technologies such as smartphone videotelephony may include reduced infection risks, increased control of disease during epidemic conditions, improved access to care, reduced stress and exposure to other pathogens during illness for better recovery, reduced time and labor costs, efficient more accessible matching of patients with particular symptoms and clinicians who are experts for such, and reduced travel while disadvantages may include privacy breaches (e.g. due to software backdoors and vulnerabilities or sale of data), dependability on Internet access and, depending on various factors, increased health care use.
Theoretically, the whole health system could benefit from telehealth. There are indications telehealth consumes fewer resources and requires fewer people to operate it with shorter training periods to implement initiatives. Commenters suggested that lawmakers may fear that making telehealth widely accessible, without any other measures, would lead to patients using unnecessary health care services. Telemedicine could also be used for connected networks between health care professionals.
Telemedicine also can eliminate the possible transmission of infectious diseases or parasites between patients and medical staff. This is particularly an issue where MRSA is a concern. Additionally, some patients who feel uncomfortable in a doctors office may do better remotely. For example, white coat syndrome may be avoided. Patients who are home-bound and would otherwise require an ambulance to move them to a clinic are also a consideration.
However, whether or not the standard of health care quality is increasing is debatable, with some literature refuting such claims. Research has reported that clinicians find the process difficult and complex to deal with. Furthermore, there are concerns around informed consent, legality issues as well as legislative issues. A recent study also highlighted that the swift and large-scale implementation of telehealth across the United Kingdom NHS Allied Health Professional (AHP) services might increase disparities in health care access for vulnerable populations with limited digital literacy. Although health care may become affordable with the help of technology, whether or not this care will be "good" is the issue. Many studies indicate high satisfaction with telemedicine among patients. Among the factors associated with a good trust in telemedicine, the use of known and user-friendly video services and confidence in the data protection policies were the two variables contributing most to trust in telemedicine.
Major problems with increasing adoption include technically challenged staff, resistance to change or habits and age of patient. Focused policy could eliminate several barriers.
A review lists a number of potentially good practices and pitfalls, recommending the use of "virtual handshakes" for confirming identity, taking consent for conducting remote consultation over a conventional meeting, and professional standardized norms for protecting patient privacy and confidentiality. It also found that the COVID-19 pandemic substantially increased, voluntarily, the adoption of telephone or video consultation and suggests that telemedicine technology "is a key factor in delivery of health care in the future".
Economic evaluations
Due to its digital nature it is often assumed that telehealth saves the health system money. However, the evidence to support this is varied. When conducting economic evaluations of telehealth services, the individuals evaluating them need to be aware of potential outcomes and extraclinical benefits of the telehealth service. Economic viability relies on the funding model within the country being examined (public vs private), the consumers willingness-to-pay, and the expected remuneration by the clinicians or commercial entities providing the services (examples of research on these topics from teledermoscopy in Australia)
In a UK telehealth trial done in 2011, it was reported that the cost of health could be dramatically reduced with the use of telehealth monitoring. The usual cost of in vitro fertilisation (IVF) per cycle would be around $15,000; with telehealth it was reduced to $800 per patient. In Alaska the Federal Health Care Access Network, which connects 3,000 healthcare providers to communities, engaged in 160,000 telehealth consultations from 2001 and saved the state $8.5 million in travel costs for just Medicaid patients.
Digital interventions for mental health conditions seem to be cost-effective compared to no intervention or non-therapeutic responses such as monitoring. However, when compared to in-person therapy or medication their added value is currently uncertain.
Beneficial enablements
Telemedicine can be beneficial to patients in isolated communities and remote regions, who can receive care from doctors or specialists far away without the patient having to travel to visit them. Recent developments in mobile collaboration technology can allow healthcare professionals in multiple locations to share information and discuss patient issues as if they were in the same place. Remote patient monitoring through mobile technology can reduce the need for outpatient visits and enable remote prescription verification and drug administration oversight, potentially significantly reducing the overall cost of medical care. It may also be preferable for patients with limited mobility, for example, patients with Parkinson's disease. Telemedicine can also facilitate medical education by allowing workers to observe experts in their fields and share best practices more easily.
During war and disasters
Remote surgery and types of videoconferencing for sharing expertise (e.g. ad hoc assistance) have been and could be used to support doctors in Ukraine during the 2022 Russian invasion of Ukraine.
Nonclinical uses
Distance education including continuing medical education, grand rounds, and patient education
administrative uses including meetings among telehealth networks, supervision, and presentations
research on telehealth
online information and health data management
healthcare system integration
asset identification, listing, and patient to asset matching, and movement
overall healthcare system management
patient movement and remote admission
Physical distancing to prevent transmission of communicable diseases
Limitations and restrictions
While many branches of medicine have wanted to fully embrace telehealth for a long time, there are certain risks and barriers which bar the full amalgamation of telehealth into best practice. For a start, it is dubious as to whether a practitioner can fully leave the "hands-on" experience behind. Although it is predicted that telehealth will replace many consultations and other health interactions, it cannot yet fully replace a physical examination, this is particularly so in diagnostics, rehabilitation or mental health. To minimise safety issues, researchers have suggested not offering remote consultations for some conditions (breathing problems, new psychosis, or acute chest pain, for example), when a parent is very concerned about a child, when a condition has not resolved as expected or has worsened, or to people who might struggle to understand or be understood (such as those with limited English or learning difficulties).
The benefits posed by telehealth challenge the normative means of healthcare delivery set in both legislation and practice. Therefore, the growing prominence of telehealth is starting to underscore the need for updated regulations, guidelines and legislation which reflect the current and future trends of healthcare practices. Telehealth enables timely and flexible care to patients wherever they may be; although this is a benefit, it also poses threats to privacy, safety , medical licensing and reimbursement. When a clinician and patient are in different locations, it is difficult to determine which laws apply to the context. Once healthcare crosses borders different state bodies are involved in order to regulate and maintain the level of care that is warranted to the patient or telehealth consumer. As it stands, telehealth is complex with many grey areas when put into practice especially as it crosses borders. This effectively limits the potential benefits of telehealth.
An example of these limitations include the current American reimbursement infrastructure, where Medicare will reimburse for telehealth services only when a patient is living in an area where specialists are in shortage, or in particular rural counties. The area is defined by whether it is a medical facility as opposed to a patient's' home. The site that the practitioner is in, however, is unrestricted. Medicare will only reimburse live video (synchronous) type services, not store-and-forward, mhealth or remote patient monitoring (if it does not involve live-video). Some insurers currently will reimburse telehealth, but not all yet. So providers and patients must go to the extra effort of finding the correct insurers before continuing. Again in America, states generally tend to require that clinicians are licensed to practice in the surgery' state, therefore they can only provide their service if licensed in an area that they do not live in themselves.
More specific and widely reaching laws, legislations and regulations will have to evolve with the technology. They will have to be fully agreed upon, for example, will all clinicians need full licensing in every community they provide telehealth services too, or could there be a limited use telehealth licence? Would the limited use licence cover all potential telehealth interventions, or only some? Who would be responsible if an emergency was occurring and the practitioner could not provide immediate help – would someone else have to be in the room with the patient at all consult times? Which state, city or country would the law apply in when a breach or malpractice occurred?
A major legal action prompt in telehealth thus far has been issues surrounding online prescribing and whether an appropriate clinician-patient relationship can be established online to make prescribing safe, making this an area that requires particular scrutiny. It may be required that the practitioner and patient involved must meet in person at least once before online prescribing can occur, or that at least a live-video conference must occur, not just impersonal questionnaires or surveys to determine need.
Telehealth has some potential for facilitating self-management techniques in health care, but for patients to benefit from it, the appropriate contact with, and relationship, between doctor and patient must be established first. This would start with an online consultation, providing patients with techniques and tools that help them participate in healthy behaviors, and initiating a collaborative partnership between health care professionals and patient. Self-management strategies fall into a broader category called patient activation, which is defined as a "patients' willingness and ability to take independent actions to manage their health." It can be achieved by increasing patients' knowledge and confidence in coping with and managing their own disease through a "regular assessment of progress [...] and problem-solving support." Teaching patients about their conditions and ways to cope with chronic illnesses will allow them to be knowledgeable about their disease and willing to manage it, improving their everyday life. Without a focus on the doctor-patient relationship and on the patient's understanding, telehealth cannot improve the quality of life of patients, despite the benefit of allowing them to do their medical check-ups from the comfort of their home.
The downsides of telemedicine include the cost of telecommunication and data management equipment and of technical training for medical personnel who will employ it. Virtual medical treatment also entails potentially decreased human interaction between medical professionals and patients, an increased risk of error when medical services are delivered in the absence of a registered professional, and an increased risk that protected health information may be compromised through electronic storage and transmission. There is also a concern that telemedicine may actually decrease time efficiency due to the difficulties of assessing and treating patients through virtual interactions; for example, it has been estimated that a teledermatology consultation can take up to thirty minutes, whereas fifteen minutes is typical for a traditional consultation. Additionally, potentially poor quality of transmitted records, such as images or patient progress reports, and decreased access to relevant clinical information are quality assurance risks that can compromise the quality and continuity of patient care for the reporting doctor. Other obstacles to the implementation of telemedicine include unclear legal regulation for some telemedical practices and difficulty claiming reimbursement from insurers or government programs in some fields. Some medical organizations have delivered position statement on the correct use of telemedicine in their field.
Another disadvantage of telemedicine is the inability to start treatment immediately. For example, a patient with a bacterial infection might be given an antibiotic hypodermic injection in the clinic, and observed for any reaction, before that antibiotic is prescribed in pill form.
Equitability is also a concern. Many families and individuals in the United States, and other countries, do not have internet access in their homes or the proper electronic devices to access services such as a laptop or smartphone.
Ethical issues
Informed consent is another issue. When telehealth includes the possibility for technical problems such as transmission errors, security breaches, or storage issues, it can impact the system's ability to communicate. It may be wise to obtain informed consent in person first, as well as having backup options for when technical issues occur. In person, a patient can see who is involved in their care (namely themselves and their clinician in a consult), but online there will be other involved such as the technology providers, therefore consent may need to involve disclosure of anyone involved in the transmission of the information and the security that will keep their information private, and any legal malpractice cases may need to involve all of those involved as opposed to what would usually just be the practitioner.
State of the market
The rate of adoption of telehealth services in any jurisdiction is frequently influenced by factors such as the adequacy and cost of existing conventional health services in meeting patient needs; the policies of governments and/or insurers with respect to coverage and payment for telehealth services; and medical licensing requirements that may inhibit or deter the provision of telehealth second opinions or primary consultations by physicians.
Projections for the growth of the telehealth market are optimistic, and much of this optimism is predicated upon the increasing demand for remote medical care. According to a recent survey, nearly three-quarters of U.S. consumers say they would use telehealth. At present, several major companies along with a bevvy of startups are working to develop a leading presence in the field.
In the UK, the Government's Care Services minister, Paul Burstow, has stated that telehealth and telecare would be extended over the next five years (2012–2017) to reach three million people.
United States
In the United States, telemedicine companies are collaborating with health insurers and other telemedicine providers to expand marketshare and patient access to telemedicine consultations.
, 95% of employers believe their organizations will continue to provide health care benefits over the next five years.
The COVID-19 pandemic drove increased usage of telehealth services in the U.S. The U.S. Centers for Disease Control and Prevention reported a 154% increase in telehealth visits during the last week of March 2020, compared to the same dates in 2019.
Switzerland
From 1999 to 2018, the University Hospital of Zurich (USZ) offered clinical telemedicine and online medical advice on the Internet. A team of doctors answered around 2500 anonymous inquiries annually, usually within 24 to 48 hours. The team consisted of up to six physicians who are specialists in clinical telemedicine at the USZ and have many years of experience, particularly in internal and general medicine. In the entire period, 59360 inquiries were sent and answered. The majority of the users were female and on average 38 years old. However, in the course of time, considerably more men and older people began to use the service. The diversity of medical queries covered all categories of the International Statistical Classification of Diseases and Related Health Problems (ICD) and correlated with the statistical frequency of diseases in hospitals in Switzerland. Most of the inquiries concerned unclassified symptoms and signs, services related to reproduction, respiratory diseases, skin diseases, health services, diseases of the eye and nervous systems, injuries and disorders of the female genital tract. As with the Swedish online medical advice service, one-sixth of the requests related to often shameful and stigmatised diseases of the genitals, gastrointestinal tract, sexually transmitted infections, obesity and mental disorders. By providing an anonymous space where users can talk about (shameful) diseases, online telemedical services empower patients and their health literacy is enhanced by providing individualized health information. The Clinical Telemedicine and Online Counselling service of the University Hospital of Zurich is currently being revised and will be offered in a new form in the future.
Developing countries
For developing countries, telemedicine and eHealth can be the only means of healthcare provision in remote areas. For example, the difficult financial situation in many African states and lack of trained health professionals has meant that the majority of the people in sub-Saharan Africa are badly disadvantaged in medical care, and in remote areas with low population density, direct healthcare provision is often very poor However, provision of telemedicine and eHealth from urban centers or from other countries is hampered by the lack of communications infrastructure, with no landline phone or broadband internet connection, little or no mobile connectivity, and often not even a reliable electricity supply.
Telemedicine in developing countries
Telemedicine in India
India has broad rural-urban population and rural India is bereaved from medical facilities, giving telemedicine a space for growth in India. Deprived education and medical professionals in rural areas is the reason behind government's ideology to use technology to bridge this gap. Remote areas not only present a number of challenges for the service providers but also for the families who are accessing these services. Since 2018, telemedicine has expanded in India. It has undertaken a new way for doctor consultations. On 25 March 2020, in the wake of COVID-19 pandemic, the Ministry of Health and Family Welfare issued India's Telemedicine Practice Guidelines. The Board of Governors entasked by the Health Ministry published an amendment to the Indian Medical Council (Professional Conduct, Etiquette and Ethics) Regulations, 2002 that gave much-needed statutory support for the practice of telemedicine in India. This sector is at an ever-growing stage with high scope of development. In April 2020, the union health ministry launched the eSanjeevani telemedicine service that operates at two levels: the doctor-to-doctor telemedicine platform, and the doctor-to-patient platform. This service crossed five million tele-consultations within a year of its launch indicating conducive environment for acceptability and growth of telemedicine in India.
Telemedicine in Sub-Saharan Africa
Sub-Saharan Africa is marked by the massive introduction of new technologies and internet access. Urban areas are facing a rapid change and development, and access to internet and health is rapidly improving. Population in remote areas however, still lack access to healthcare and modern technologies. Some people in rural regions must travel more between 2 and 6 hours to reach the closest healthcare facilities of their country. leaving room for telehealth to grow and reach isolated people in the near future.
Internet via satellite in rural areas
The Satellite African eHEalth vaLidation (SAHEL) demonstration project has shown how satellite broadband technology can be used to establish telemedicine in such areas. SAHEL was started in 2010 in Kenya and Senegal, providing self-contained, solar-powered internet terminals to rural villages for use by community nurses for collaboration with distant health centers for training, diagnosis and advice on local health issues. Those methods can have major impact on both health professionals to get and provide training from remote areas, and on the local population who can receive care without traveling long distances. Some non-profits provide internet to rural places around the world using a mobile VSAT terminal. This VSAT terminal equips remote regions allowing them to alert the world when there is a medical emergency, resulting in a rapid deployment or response from developed countries. Technologies such as the ones used by MAF allows health professionals in remote clinics to have internet access, making consultations much easier, both for patients and doctors.
In 2014, the government of Luxembourg, along with satellite operators and NGOs established SATMED, a multilayer eHealth platform to improve public health in remote areas of emerging and developing countries, using the Emergency.lu disaster relief satellite platform and the Astra 2G TV satellite. SATMED was first deployed in response to a report in 2014 by German Doctors of poor communications in Sierra Leone hampering the fight against Ebola, and SATMED equipment arrived in the Serabu clinic in Sierra Leone in December 2014. In June 2015 SATMED was deployed at Maternité Hospital in Ahozonnoude, Benin to provide remote consultation and monitoring, and is the only effective communication link between Ahozonnoude, the capital and a third hospital in Allada, since land routes are often inaccessible due to flooding during the rainy season.
History
The development and history of telehealth or telemedicine (terms used interchangeably in literature) is deeply rooted in the history and development in not only technology but also society itself. Humans have long sought to relay important messages through torches, optical telegraphy, electroscopes, and wireless transmission. Early forms of telemedicine achieved with telephone and radio have been supplemented with videotelephony, advanced diagnostic methods supported by distributed client/server applications, and additionally with telemedical devices to support in-home care.
In the 21st century, with the advent of the internet, portable devices and other such digital devices are taking a transformative role in healthcare and its delivery.
Earliest instances
Although, traditional medicine relies on in-person care, the need and want for remote care has existed from the Roman and pre-Hippocratic periods in antiquity. The elderly and infirm who could not visit temples for medical care sent representatives to convey information on symptoms and bring home a diagnosis as well as treatment. In Africa, villagers would use smoke signals to warn neighboring villages of disease outbreak. The beginnings of telehealth have existed through primitive forms of communication and technology. The exact date of origin for Telehealth is unknown, but it was known to have been used during the Bubonic Plague. That version of telehealth was far different from how we know it today. During that time, they were communicating by heliograph and bonfire. Those were used to notify other groups of people about famine and war. Those are not using any form of technology yet but are starting to spread the idea of connectivity among groups of people who geographically could not be together.
1800s to early 1900s
As technology developed and wired communication became increasingly commonplace, the ideas surrounding telehealth began emerging. The earliest telehealth encounter can be traced to Alexander Graham Bell in 1876, when he used his early telephone as a means of getting help from his assistant Mr. Watson after he spilt acid on his trousers. Another instance of early telehealth, specifically telemedicine was reported in The Lancet in 1879. An anonymous writer described a case where a doctor successfully diagnosed a child over the telephone in the middle of the night. This Lancet issue, also further discussed the potential of Remote Patient Care in order to avoid unnecessary house visits, which were part of routine health care during the 1800s. Other instances of telehealth during this period came from the American Civil War, during which telegraphs were used to deliver casualty/mortality lists, medical care to soldiers, and ordering further medical supplies.
As the 1900s started, physicians quickly found a use for the telephone making it a prime communication channel to contact patients and other physicians. Over the next fifty-plus years, the telephone was a staple for medical communication. As the 1930s came around, radio communication played a key role, especially during World War I. It was specifically used to communicate with remote areas such as Alaska and Australia. They used the radio to communicate medical information. During the Vietnam War, radio communication had become more advanced and was now used to send medical teams in helicopters to help. This then brought together the Aerial Medical Service (AMS) who used telegraphs, radios, and planes to help care for people who lived in remote areas.
From the late 1800s to the early 1900s the early foundations of wireless communication were laid down. Radios provided an easier and near instantaneous form of communication. The use of radio to deliver healthcare became accepted for remote areas. The Royal Flying Doctor Service of Australia is an example of the early adoption of radios in telehealth.
In 1925 the inventor Hugo Gernsback wrote an article for the magazine Science and Invention which included a prediction of a future where patients could be treated remotely by doctors through a device he called a "teledactyl". His descriptions of the device are similar to what would later become possible with new technology.
Mid-1900s to 1980s
When the American National Aeronautics and Space Administration (NASA) began plans to send astronauts into space, the need for telemedicine became clear. In order to monitor their astronauts in space, telemedicine capabilities were built into the spacecraft as well as the first spacesuits. Additionally, during this period, telehealth and telemedicine were promoted in different countries especially the United States and Canada. Carrier Sekani Family Services helped pioneer telehealth in British Columbia and Canada, according to its CEO Warner Adam. After the telegraph and telephone started to successfully help physicians treat patients from remote areas, telehealth became more recognized. Technological advancements occurred when NASA sent men to space. Engineers for NASA created biomedical telemetry and telecommunications systems. NASA technology monitored vitals such as blood pressure, heart rate, respiration rate, and temperature. After the technology was created, it then became the base of telehealth medicine for the public.
Massachusetts General Hospital and Boston's Logan International Airport had a role in the early use of telemedicine, which more or less coincided with NASA's foray into telemedicine through the use of physiologic monitors for astronauts. On October 26, 1960, a plane struck a flock of birds upon takeoff, killing many passengers and leaving a number wounded. Due to the extreme complexity of trying to get all the medical personnel out from the hospital, the practical solution became telehealth. This was expanded upon in 1967, when Kenneth Bird at Massachusetts General founded one of the first telemedicine clinics. The clinic addressed the fundamental problem of delivering occupational and emergency health services to employees and travellers at the airport, located three congested miles from the hospital. Clinicians at the hospital would provide consultation services to patients who were at the airport. Consultations were achieved through microwave audio as well as video links. The airport began seeing over a hundred patients a day at its nurse-run clinic that cared for victims of plane crashes and other accidents, taking vital signs, electrocardiograms, and video images that were sent to Massachusetts General. Over 1,000 patients are documented as having received remote treatment from doctors at MGH using the clinic's two-way audiovisual microwave circuit. One notable story featured a woman who got off a flight in Boston and was experiencing chest pain. They performed a workup at the airport, took her to the telehealth suite where Raymond Murphy appeared on the television, and had a conversation with her. While this was happening, another doctor took notes and the nurses took vitals and any test that Murphy ordered. At this point, telehealth was becoming more mainstream and was starting to become more technologically advanced, which created a viable option for patients.
In 1964, the Nebraska Psychiatric Institute began using television links to form two-way communication with the Norfolk State Hospital which was 112 miles away for the education and consultation purposes between clinicians in the two locations.
In 1972 the Department of Health, Education and Welfare in the United States approved funding for seven telemedicine projects across different states. This funding was renewed and two further projects were funded the following year.
In March 1972, the San Bernardino County Medical Society officially implemented its Tel-Med program, a system of prerecorded health-related messages, with a log of 50 tapes. The nonprofit initiative began in 1971 as a local medical project to ease the doctor shortage in the expanding San Bernardino Valley and improve the public's access to sound medical information. It covered subjects ranging from cannabis to vaginitis. In January 1973, in response to the developing "London flu" epidemic hitting California and the country, a tape providing information on the disease was on air within a week after news broke of the flu spreading in the state. That spring, programs were implemented in San Diego and Indianapolis, Indiana, signaling a national acceptance of the concept. By 1979, its system offered messages on over 300 different subjects, 200 of which were available in Spanish as well as English, and serviced over 65 million people in 180 cities around the country.
1980s to 1990s – maturation and renaissance
Telehealth projects underway before and during the 1980s would take off but fail to enter mainstream healthcare. As a result, this period of telehealth history is called the "maturation" stage and made way for sustainable growth. Although state funding in North America was beginning to run low, different hospitals began to launch their own telehealth initiatives. NASA provided an ATS-3 satellite to enable medical care communications of American Red Cross and Pan American Health Organization response teams, following the 1985 Mexico City earthquake. The agency then launched its SateLife/HealthNet programme to increase health service connectivity in developing countries. In 1997, NASA sponsored Yale's Medical Informatics and Technology Applications Consortium project.
Florida first experimented with "primitive" telehealth in its prisons during the latter 1980s. Working with Doctors Oscar W. Boultinghouse and Michael J. Davis, from the early 1990s to 2007; Glenn G. Hammack led the University of Texas Medical Branch (UTMB) development of a pioneering telehealth program in Texas state prisons. The three UTMB alumni would, in 2007, co-found telehealth provider NuPhysician.
The first interactive telemedicine system, operating over standard telephone lines, designed to remotely diagnose and treat patients requiring cardiac resuscitation (defibrillation) was developed and launched by an American company, MedPhone Corporation, in 1989. A year later under the leadership of its President/CEO S Eric Wachtel, MedPhone introduced a mobile cellular version, the MDPhone. Twelve hospitals in the U.S. served as receiving and treatment centers.
At-home virtual care
As the expansion of telehealth continued in 1990 Maritime Health Services (MHS) was a big part of the initiation for occupational health services. They sent a medical officer aboard the Pacific trawler that allowed for round-the-clock communication with a physician. The system that allows for this is called the Medical Consultation Network (MedNet). MedNet is a video chatting system that has live audio and visual so the physician on the other end of the call can see and hear what is happening. MetNet can be used from anywhere, not just aboard ships. Being able to provide onsite visual information allows remote patients expert emergency help and medical attention that saves money as well as lives. This has created a demand for at-home monitoring. At-home care has also become a large part of telehealth. Doctors or nurses will now give pre-op and post-op phone calls to check-in. There are also companies such as Lifeline, which give the elderly a button to press in case of an emergency. That button will automatically call for emergency help. If someone has surgery and then is sent home, telehealth allows physicians to see how the patient is progressing without them having to stay in the hospital. TeleDiagnostic Systems of San Francisco is a company that has created a device that monitors sleep patterns, so people with sleep disorders do not have to stay the night at the hospital. Another at-home device that was created was the Wanderer, which was attached to Alzheimer's patients or people who had dementia. It was attached to them so when they wandered off it notified the staff to allow them to go after them. All these devices allowed healthcare beyond hospitals to improve, which means that more people are being helped efficiently.
2000s to present
The advent of high-speed Internet, and the increasing adoption of ICT in traditional methods of care, spurred advances in telehealth delivery. Increased access to portable devices, like laptops and mobile phones, made telehealth more plausible; the industry then expanded into health promotion, prevention and education.
In 2002, G. Byron Brooks, a former NASA surgeon and engineer who had also helped manage the UTMB Telemedicine program, co-founded Teladoc in Dallas, Texas, which was then launched in 2005 as the first national telehealth provider.
In the 2010s, integration of smart home telehealth technologies, such as health and wellness devices, software, and integrated IoT, has accelerated the industry. Healthcare organizations are increasingly adopting the use of self-tracking and cloud-based technologies, and innovative data analytic approaches to accelerate telehealth delivery.
In 2015, Mercy Health system opened Mercy Virtual, in Chesterfield, Missouri, the world's first medical facility dedicated solely to telemedicine.
COVID-19
Telehealth expanded significantly during the COVID-19 pandemic, becoming a vital means of medical communication. It allows doctors to return to humanizing the patient. It forces them to listen to what people have to say and from there make a diagnosis. Studies have demonstrated high trust in telehealth expressed by patients during the COVID-19 pandemic. Some researchers claim this creates an environment that encourages greater vulnerability among patients in self disclosure in the practice of narrative medicine. Telehealth allows for Zoom calls and video chats from across the world checking in on patients and speaking to physicians. Universities are now ensuring that medical students graduate with proficient telehealth communication skills. Experts suggest that telehealth has become a vital part of medical care; with more virtual options becoming available. The pandemic era also identified the "potential to significantly improve global health equity" through telehealth and other virtual care technologies.
| Biology and health sciences | Diagnostics | Health |
1178458 | https://en.wikipedia.org/wiki/Flickr | Flickr | Flickr ( ; ) is an image hosting and video hosting service, as well as an online community, founded in Canada and headquartered in the United States. It was created by Ludicorp in 2004 and was previously a common way for amateur and professional photographers to host high-resolution photos. It has changed ownership several times and has been owned by SmugMug since April 20, 2018.
Flickr had a total of 112 million registered members and more than 3.5 million new images uploaded daily. On August 5, 2011, the site reported that it was hosting more than 6 billion images. In 2024, it was reported as having shared 10 billion photos and accepting 25 million per day.
Photos and videos can be accessed from Flickr without the need to register an account, but an account must be made to upload content to the site. Registering an account also allows users to create a profile page containing photos and videos that the user has uploaded and also grants the ability to add another Flickr user as a contact. For mobile users, Flickr has official mobile apps for iOS, Android, and an optimized mobile site.
History
Flickr was launched on February 10, 2004, by Ludicorp, a Vancouver-based company founded by Stewart Butterfield and Caterina Fake. The service emerged from tools originally created for Ludicorp's Game Neverending, a web-based massively multiplayer online game. Flickr proved a more feasible project, and ultimately Game Neverending was shelved. Butterfield later launched a similar online game, Glitch, which was shut down on November 14, 2012.
Early versions of Flickr focused on a chat room called FlickrLive, with real-time photo exchange capabilities. The successive evolutions focused more on the uploading and filing back-end for individual users and the chat room was buried in the site map. It was eventually dropped as Flickr's back-end systems evolved away from Game Neverendings codebase. Key features of Flickr not initially present are tags, marking photos as favorites, group photo pools and interestingness, for which a patent was granted.
In addition to being a popular website for users to share and embed personal photographs and an online community, in 2004, the service was widely used by photo researchers and by bloggers to host images that they embed in blogs and social media.
Yahoo! acquired Ludicorp and Flickr on March 20, 2005. The acquisition reportedly cost between $22 million and $25 million. During the week of June 26, 2005 to July 2, 2005, all content was migrated from servers in Canada to servers in the United States, and all resulting data became subject to United States federal law. On May 3, 2007, Yahoo! announced that Yahoo! Photos would close down on September 20, 2007, after which all photos would be deleted; users were encouraged to migrate to Flickr. On January 31, 2007, Flickr announced that, to continue using the service, "Old Skool" members (those who had joined before the Yahoo! acquisition) would be required to associate their account with a Yahoo! identity by March 15, 2007. That move was criticized by some users.
Flickr upgraded its services from "beta" to "gamma" status on May 16, 2006, the changes attracted positive attention from Lifehacker. On December 13, 2006, upload limits on free accounts were increased to 100 MB a month (from 20 MB) and were removed from Flickr Pro accounts, which originally had a 2 GB per month limit. On April 9, 2008, Flickr began allowing paid subscribers to upload videos, limited to 90 seconds in length and 150 MB in size. On March 2, 2009, Flickr added the facility to upload and view HD videos, and began allowing free users to upload normal-resolution video. At the same time, the set limit for free accounts was lifted. In 2009, Flickr announced a partnership with Getty Images in which selected users could submit photographs for stock photography usage and receive payment. On June 16, 2010, this was changed so that users could label images as suitable for stock use themselves.
On May 20, 2013, Flickr launched the first stage of a major site redesign, introducing a "Justified View" close-spaced photo layout browsed via "infinite scrolling" and adding new features, including one terabyte of free storage for all users, a scrolling home page (mainly of contacts photos and comments) and updated Android app. The Justified View is paginated between 72 and 360 photos per page but unpaginated in search result presentation. Tech Radar described the new style Flickr as representing a "sea change" in its purpose. Many users criticized the changes, and the site's help forum received thousands of negative comments. On March 25, 2014, Flickr's New Photo Experience, a user interface redesign, left beta.
On May 7, 2015, Yahoo! overhauled the site, adding a revamped Camera Roll, a new way to upload photos, and upgraded the site's apps. The new Uploadr application was made available for Macs, Windows and mobile devices.
In 2018, Yahoo! (owned by Verizon at that point) sold Flickr to SmugMug.
In early May 2019, SmugMug announced the migration of Flickr data, involving 100+ million accounts and billions of photos and videos, from the servers of former owner Yahoo! to Amazon Web Services (AWS), in a planned 12-hour transition to occur on May 22, 2019.
In May 2023, Flickr announced the development of the Print Shop feature that was being tested with a list of approved sellers. The Print shop feature allows photographers to sell prints via a storefront, and allows purchases from consumers. The feature was to allow only approved members access to it, but the criteria for that were yet to be announced.
Corporate changes
On June 13, 2008, Flickr co-founder Stewart Butterfield announced his resignation on July 12, 2008, which followed that of his wife and co-founder Caterina Fake, who left the company on the same day. Butterfield wrote a humorous resignation letter to Brad Garlinghouse.
On December 14, 2008, The Guardian reported that three employees had been laid off as Yahoo! continued to reduce its workforce and, on November 30, 2010, CNET reported Yahoo! was on the verge of a major layoff, affecting 10% to 20% of its workforce. Flickr was specifically named as a target for these layoffs.
On June 13, 2017, Verizon Communications acquired Yahoo!, including Flickr. Verizon reorganized Yahoo!, along with AOL, into a new umbrella company, Oath, which was renamed as Verizon Media on January 8, 2019.
On April 20, 2018, SmugMug acquired Flickr from Verizon's Oath and put an end to Flickr 1 TB storage plan for free users. Those users had until February 5, 2019, to convert to "Pro" accounts or their photo streams would be reduced to a maximum of 1,000 pictures. The deadline was later extended to March 12, 2019. The reasons cited were that the existing model was unsustainable by a medium-sized company which could not get revenues by selling profiles of the users. The sentiment was generally agreed on among the professionals. This policy, however, was never implemented and was abandoned in March, 2022 in favor of a policy that restricted content unless the user upgraded and paid for a Pro account.
Features
Accounts
Flickr has always offered two types of accounts: free and paid. Until January 7, 2019, free accounts had up to 1 TB of storage. On January 8, 2019, the account offerings changed. The free option is limited to 1,000 photos or videos stored, with videos limited to three minutes. After January 8, 2019, members over the limit could no longer upload new photos to Flickr. On February 5, 2019, a free account's older content would be deleted automatically if it contains more than 1,000 photos and they do not subscribe to the paid service tier, with the exception of content that was already uploaded with a Creative Commons copyright license before November 1, 2018.
The paid option features "unlimited" storage, advanced statistics, advertising-free browsing, videos up to 10 minutes in length, "premier" customer service, and promotional offers with other partners.
In May 2011, Flickr added an option to easily reverse an account termination, motivated by the accidental deletion of a Flickr user's account, and public reporting of its protracted restoration. Flickr may delete accounts without giving any reason or warning to the account's owner.
As a result of the SmugMug buyout, Flickr added the ability for users to download all of their account data, such as photo albums, contacts, and comments.
Organization
The images a photographer uploads to Flickr go into their sequential "photostream", the basis of a Flickr account. All photostreams can be displayed as a justified view, a slide show, a "detail" view or a date stamped archive. Clicking on a photostream image opens it in the interactive "photopage" alongside data, comments and facilities for embedding images on external sites.
Users may label their uploaded images with titles and descriptions, and images may be tagged, either by the uploader or by other users, if the uploader permits it. These text components enable computer searching of Flickr. Flickr was an early website to implement tag clouds, which were used until 2013, providing access to images tagged with the most popular keywords. Tagging was further revised in the photopage redesign of March 2014. Flickr has been cited as a prime example of effective use of folksonomy.
Users can organize their Flickr photos into "albums" (formerly "sets") which are more flexible than the traditional folder-based method of organizing files, as one photo can belong to one album, many albums, or none at all. Flickr provides code to embed albums into blogs, websites and forums. Flickr albums represent a form of categorical metadata rather than a physical hierarchy. Geotagging can be applied to photos in albums, and any albums with geotagging can be related to a map using imapflickr. The resulting map can be embedded in a website. Flickr albums may be organized into "collections", which can themselves be further organized into higher-order collections.
Organizr is a Web application for organizing photos within a Flickr account that can be accessed through the Flickr interface. It allows users to modify tags, descriptions and set groupings, and to place photos on a world map (a feature provided in conjunction with Yahoo! Maps). It uses Ajax to emulate the look, feel and quick functionality of desktop-based photo-management applications, such as Google's Picasa and F-Spot. Users can select and apply changes to multiple photos at a time, as an alternative to the standard Flickr interface for editing.
Access control
Flickr provides both private and public image storage. A user uploading an image can set privacy controls that determine who can view the image. A photo can be flagged as either public or private. Private images are visible by default only to the uploader, but they can also be marked as viewable by friends or family. Privacy settings also can be decided by adding photographs from a user's photostream to a "group pool". If a group is private all the members of that group can see the photo. If a group is public the photo becomes public as well. Flickr also provides a "contact list" which can be used to control image access for a specific set of users in a way similar to that of LiveJournal. In November 2006, Flickr created a "guest pass" system that allows private photos to be shared with non-Flickr members. This setting allows sets or all photos under a certain privacy category (friends or family) to be shared. Many members allow their photos to be viewed by anyone, forming a large collaborative database of categorized photos. By default, other members can leave comments about any image they have permission to view and, in many cases, can add to the list of tags associated with an image.
Interaction and compatibility
The core functionality of the site relies on standard HTML and HTTP features, allowing for wide compatibility among platforms and browsers; Flickr's functionality includes RSS and Atom feeds and an API that enables independent programmers to expand its services. This includes a large number of third-party Greasemonkey scripts that enhance and extend the functionality of Flickr. In 2006, Flickr was the second most extended site on userscripts.org. Organizr and most of Flickr's other text-editing and tagging interfaces use Ajax, with which most modern browsers are compliant. Images can be posted to the user's photostream via email attachments, which enables direct uploads from many smartphones and applications. Flickr uses the Geo microformat on over three million geotagged images.
According to the company, Flickr is hosted on 62 databases across 124 servers, with about 800,000 user accounts per pair of servers. Based on information compiled by highscalability.com, the MySQL databases are hosted on servers that are Linux-based (from Red Hat), with a software platform that includes Apache, PHP (with PEAR and Smarty), shards, Memcached, Squid, Perl, ImageMagick and Java; the system administration tools include Ganglia, SystemImager, Subcon, and CVSup.
Signed-in Flickr users can "Follow" the Photostreams of other Flickr photographers. Reciprocating this process is optional. A user's homepage contains a stream of their Contacts' photos at 2/3 screen size.
Groups are another major means of interaction with fellow members of Flickr around common photography interests. A Flickr Group can be started by any Flickr user, who becomes its administrator and can appoint moderators. Groups may either be open access or invitation-only, and most have an associated pool of photos. The administrator of the Flickr group can monitor and set restrictions for the group, assign awards to members, and may curate and organize the photo content. Recent uploads to a group will sometimes appear on its members' homepages. Group photo pools may be displayed in the "Justified View" or as a slideshow.
"Galleries" of photos from other photostreams may be curated by any signed-up Flickr user, provided the feature is not disabled by the photo's uploader, these are then publicly viewable.
Any Flickr user can post comments to a Flickr photo on its photopage, unless this has been disabled by the uploader, and users can "favorite" a photo. A user's favorites can be viewed in a justified or slideshow display.
Users of Windows Photo Gallery, Apple's iPhoto (version 8), Adobe's Lightroom 3.2, Apple's Aperture (version 3.0), darktable, and digiKam have the ability to upload their photos directly to Flickr. They can also automatically update their status on other social networking services when they upload their images to Flickr. Flickr provides a desktop client for Mac OS X and Windows that allows users to upload photos without using the web interface. Uploadr allows drag-and-drop batch uploading of photos, the setting of tags and descriptions for each batch, and the editing of privacy settings.
Flickr has entered into partnerships with many third parties:
Flickr had a partnership with the Picnik online photo-editing application that included a reduced-feature version of Picnik built into Flickr as a default photo editor. On April 5, 2012, Flickr replaced Picnik with Aviary as its default photo editor.
In addition to commercial mapping data, Flickr now uses OpenStreetMap mapping for various cities; this began with Beijing during the run-up to the 2008 Olympic games. , this is used for Baghdad, Beijing, Kabul, Sydney and Tokyo. OpenStreetMap data is collected by volunteers and is available under the Open Database License.
Flickr offers printing of various forms of merchandise, including business cards, photo books, stationery, personalized credit cards and large-size prints from companies such as Moo, Blurb, Tiny Prints, Capital One, Imagekind, and QOOP.
The Flickr partnership with Getty Images to sell stock photos from users is under review as of early 2014.
Filtering
In March 2007, Flickr added new content filtering controls that let members specify by default what types of images they generally upload (photo, art/illustration, or screenshot) and how "safe" (i.e., unlikely to offend others) their images are, as well as specify that information for specific images individually. Individual images are assigned to one of three categories: "safe", "moderate" and "restricted". Users can specify the same criteria when searching for images. There are some restrictions on searches for certain types of users: non-members must always use SafeSearch, which omits images noted as potentially offensive, while members whose Yahoo! accounts indicate that they are underage may use SafeSearch or moderate SafeSearch, but cannot turn SafeSearch off completely. The system achieves a fairly good separation of family-friendly photos and adult content; generic image searches normally produce no pornographic results, with the visibility of adult content restricted to users and dedicated Flickr communities who have opted into viewing it.
Flickr has used this filtering system to change the level of accessibility to "unsafe" content for entire nations, including South Korea, Hong Kong and Germany. In summer 2007, German users staged a "revolt" over being assigned the user rights of a minor.
Licensing
Flickr offers users the ability to either release their images under certain common usage licenses or label them as "all rights reserved". The licensing options primarily include the Creative Commons 2.0 attribution-based and minor content-control licenses – although jurisdiction and version-specific licenses cannot be selected. As with "tags", the site allows easy searching of only those images that fall under a specific license.
On January 16, 2008, Flickr launched a program called "The Commons on Flickr." Several international cultural institutions share images using a "no known copyright restrictions" through the program. According to Flickr, the goal of the program is to "firstly show you hidden treasures in the world's public photography archives, and secondly to show how your input and knowledge can help make these collections even richer." Participants include the National Museum of Denmark, Powerhouse Museum, George Eastman Museum, Library of Congress, Nationaal Archief, National Archives and Records Administration, National Library of Scotland, State Library of New South Wales, and Smithsonian Institution.
In May 2009, White House official photographer Pete Souza began using Flickr as a conduit for releasing White House photos. The photos were initially posted with a Creative Commons Attribution license requiring that the original photographers be credited. Flickr later created a new license which identified them as "United States Government Work", which does not carry any copyright restrictions.
In March 2015, Flickr added the Creative Commons Public Domain Mark and Creative Commons Zero (CC0) to its licensing options. The Public Domain Mark is meant for images that are no longer protected by copyright. CC0 is used for works that are still protected by copyright or neighbouring rights, but where the rightsholders choose to waive all those rights.
Reception
Flickr became an immediate success and was seen as a successful example of "Web 2.0", and a year later was purchased by Yahoo!. Initially, the site was most popular with professional photographers and graphic designers as well as bloggers who used it as an image repository. In 2007, Flickr was the 19th most popular website on the Internet according to its Alexa Rank. However, since then, its popularity has declined relative to social media platforms with photo sharing capabilities (such as Facebook and Instagram), as well as cloud file storage services (such as Dropbox). By 2021, Flickr's Alexa Rank had declined significantly, yet indicated that the website was still among the top 500 most popular websites globally.
Controversies
Censorship
On June 12, 2007, in the wake of the rollout of localized language versions of the site, Flickr implemented a user-side rating system for filtering out potentially controversial photos. Simultaneously, users with accounts registered with Yahoo! subsidiaries in Germany, Singapore, Hong Kong, and Korea were prevented from viewing photos rated "moderate" or "restricted" on the three-part scale used. Many Flickr users, particularly in Germany, protested against the new restrictions, claiming unwanted censorship from Flickr and Yahoo.
Flickr management, unwilling to go into legal details, implied that the reason for the stringent filtering was some particularly strict age-verification laws in Germany. The issue received attention in the German national media, especially in online publications. Initial reports indicated that Flickr's action was a sensible, if unattractive, precaution against prosecution, although later coverage implied that Flickr's action may have been unnecessarily strict. On June 20, 2007, Flickr reacted by granting German users access to "moderate" (but not "restricted") images, and hinted at a future solution for Germany, involving advanced age-verification procedures.
Since June 1, 2009, Flickr has been blocked in China in advance of the 20th anniversary of the Tiananmen Square protests of 1989.
Copyright enforcement
Michael Arrington of TechCrunch and the Electronic Frontier Foundation have criticised Flickr for its heavy-handed implementation of the Digital Millennium Copyright Act (DMCA) and Online Copyright Infringement Liability Limitation Act (OCILLA). Under OCILLA, a service provider such as Flickr is obliged to delete or disable access to content as soon as they receive an official notice of infringement, to maintain protection from liability. After having one of his own pictures taken down following an incorrect DMCA claim, British comedian Dave Gorman researched the issue and concluded that if the Flickr user was not based in the United States or they were, but the person filing the notice of infringement was not, Flickr deleted the disputed content immediately. Even if the user could successfully demonstrate that the content did not infringe upon any copyright, Flickr did not, according to Gorman, replace the deleted content. He argued that this was contrary to its obligations in responding to a counter-notice. Shortly afterward, Flickr changed its policy.
In 2019, Flickr added new theft detection tool options to certain users. Some subscribers will be provided "copy-protection tools that can detect if their images have been used without permission," the BBC reported in 2019, noting "Flickr Pro subscribers will be able to monitor up to 1,000 images and send automated copyright claims to people or companies that use their photos."
Sale of Creative Commons-licensed photos
In November 2014, Flickr announced that it would sell wall-sized prints of photos from the service that are licensed under Creative Commons licenses allowing commercial use. Although its use of the photos in this manner is legal and allowed under the licenses, Flickr was criticized by users for what they perceived to be unfair exploitation of artists' works, as all the profits from these offerings go to Yahoo! and are not shared with their respective photographers, and users were not given a means of opting-out from the program without placing their photos under a more restrictive non-commercial license. By contrast, a similar opt-in program for "licensed" photos does give photographers a 51% share of sales. On December 19, 2014, Flickr General Manager Bernardo Hernandez announced they would pull all Creative Commons-licensed images from the program and issue refunds, stating that "Subsequently, we'll work closely with Creative Commons to come back with programs that align better with our community values."
Deletion of files of non-paying users
On November 1, 2018, Flickr announced new restrictions for its users.
On January 8, 2019, non-paying users would only able to upload up to 1000 files free of charge.
Deletion of the oldest files, determined by the upload date, was scheduled to begin on March 12, 2019, until the limit of 1000 files was met. The size of the individual files would not be relevant. Alternatively, users could upgrade to Pro subscription for US$60 per year.
On March 17, 2022, Flickr revealed that it had not in fact deleted any photos for exceeding storage limits. However, it announced that it would soon implement a policy limiting the sharing of "moderate" or "restricted" content to Pro users, and limiting free users to 50 "non-public" images. Images beyond these limits would be at risk for deletion.
| Technology | Multimedia | null |
1178938 | https://en.wikipedia.org/wiki/Last%20Glacial%20Maximum | Last Glacial Maximum | The Last Glacial Maximum (LGM), also referred to as the Last Glacial Coldest Period, was the most recent time during the Last Glacial Period where ice sheets were at their greatest extent between 26,000 and 20,000 years ago.
Ice sheets covered much of Northern North America, Northern Europe, and Asia and profoundly affected Earth's climate by causing a major expansion of deserts, along with a large drop in sea levels.
Based on changes in position of ice sheet margins dated via terrestrial cosmogenic nuclides and radiocarbon dating, growth of ice sheets in the southern hemisphere commenced 33,000 years ago and maximum coverage has been estimated to have occurred sometime between 26,500 years ago and 20,000 years ago. After this, deglaciation caused an abrupt rise in sea level. Decline of the West Antarctica ice sheet occurred between 14,000 and 15,000 years ago, consistent with evidence for another abrupt rise in the sea level about 14,500 years ago. Glacier fluctuations around the Strait of Magellan suggest the peak in glacial surface area was constrained to between 25,200 and 23,100 years ago.
There are no agreed dates for the beginning and end of the LGM, and researchers select dates depending on their criteria and the data set consulted. Jennifer French, an archeologist specialising in the European Palaeolithic, dates its onset at 27,500 years ago, with ice sheets at their maximum by around 26,000 years ago and deglaciation commencing between 20,000 and 19,000 years ago. The LGM is referred to in Britain as the Dimlington Stadial, dated to between 31,000 and 16,000 years ago.
Glacial climate
The average global temperature about 21,000 years ago was about 6 °C (11 °F) colder than today. According to the United States Geological Survey (USGS), permanent summer ice covered about 8% of Earth's surface and 25% of the land area during the last glacial maximum. The USGS also states that sea level was about lower than in present times (2012). When comparing to the present, the average global temperature was for the 2013–2017 period. As of 2012 about 3.1% of Earth's surface and 10.7% of the land area is covered in year-round ice.
Carbon sequestration in the highly stratified and productive Southern Ocean was essential in producing the LGM. The formation of an ice sheet or ice cap requires both prolonged cold and precipitation (snow). Hence, despite having temperatures similar to those of glaciated areas in North America and Europe, East Asia remained unglaciated except at higher elevations. This difference was because the ice sheets in Europe produced extensive anticyclones above them. These anticyclones generated air masses that were so dry on reaching Siberia and Manchuria that precipitation sufficient for the formation of glaciers could never occur (except in Kamchatka where these westerly winds lifted moisture from the Sea of Japan). The relative warmth of the Pacific Ocean due to the shutting down of the Oyashio Current and the presence of large east–west mountain ranges were secondary factors that prevented the development of continental glaciation in Asia.
All over the world, climates at the Last Glacial Maximum were cooler and almost everywhere drier. In extreme cases, such as South Australia and the Sahel, rainfall could have been diminished by up to 90% compared to the present, with flora diminished to almost the same degree as in glaciated areas of Europe and North America. Even in less affected regions, rainforest cover was greatly diminished, especially in West Africa where a few refugia were surrounded by tropical grasslands. The Amazon rainforest was split into two large blocks by extensive savanna, and the tropical rainforests of Southeast Asia probably were similarly affected, with deciduous forests expanding in their place except on the east and west extremities of the Sundaland shelf. Only in Central America and the Chocó region of Colombia did tropical rainforests remain substantially intact – probably due to the extraordinarily heavy rainfall of these regions. Most of the world's deserts expanded. Exceptions were in what is the present-day Western United States, where changes in the jet stream brought heavy rain to areas that are now desert and large pluvial lakes formed, the best known being Lake Bonneville in Utah. This also occurred in Afghanistan and Iran, where a major lake formed in the Dasht-e Kavir.
In Australia, shifting sand dunes covered half the continent, while the Chaco and Pampas in South America became similarly dry. Present-day subtropical regions also lost most of their forest cover, notably in eastern Australia, the Atlantic Forest of Brazil, and southern China, where open woodland became dominant due to much drier conditions. In northern China – unglaciated despite its cold climate – a mixture of grassland and tundra prevailed, and even here, the northern limit of tree growth was at least 20° farther south than today. In the period before the LGM, many areas that became completely barren desert were wetter than they are today, notably in southern Australia, where Aboriginal occupation is believed to coincide with a wet period between 40,000 and 60,000 years Before Present (BP). In New Zealand and neighbouring regions of the Pacific, temperatures may have been further depressed during part of the LGM by the world's most recent supervolcanic eruption, the Oruanui eruption, approximately 25,500 years BP.
However, it is estimated that during the LGM, low-to-mid latitude land surfaces at low elevation cooled on average by 5.8 °C relative to their present-day temperatures, based on an analysis of noble gases dissolved in groundwater rather than examinations of species abundances that have been used in the past.
World impact
During the Last Glacial Maximum, much of the world was cold, dry, and inhospitable, with frequent storms and a dust-laden atmosphere. The dustiness of the atmosphere is a prominent feature in ice cores; dust levels were as much as 20 to 25 times greater than they are in the present. This was probably due to a number of factors: reduced vegetation, stronger global winds, and less precipitation to clear dust from the atmosphere. The massive sheets of ice locked away water, lowering the sea level, exposing continental shelves, joining land masses together, and creating extensive coastal plains. The ice sheets also changed the atmospheric circulation, causing the northern Pacific and Atlantic oceans to cool and produce more clouds, which amplified the global cooling as the clouds reflected even more sunlight. During the LGM, 21,000 years ago, the sea level was about 125 meters (about 410 feet) lower than it is today. Across most of the globe, the hydrological cycle slowed down, explaining increased aridity in many regions of the world.
Africa and the Middle East
In Africa and the Middle East, many smaller mountain glaciers formed, and the Sahara and other sandy deserts were greatly expanded in extent. The Atlantic deep sea sediment core V22-196, extracted off the coast of Senegal, shows a major southward expansion of the Sahara.
The Persian Gulf averages about 35 metres in depth and the seabed between Abu Dhabi and Qatar is even shallower, being mostly less than 15 metres deep. For thousands of years the Ur-Shatt (a confluence of the Tigris-Euphrates Rivers) provided fresh water to the Gulf, as it flowed through the Strait of Hormuz into the Gulf of Oman. Bathymetric data suggests there were two palaeo-basins in the Persian Gulf. The central basin may have approached an area of 20,000 km2, comparable at its fullest extent to lakes such as Lake Malawi in Africa. Between 12,000 and 9,000 years ago much of the Gulf's floor was not covered by water, only being flooded by the sea after 8,000 years BP.
It is estimated that annual average temperatures in Southern Africa were 6 °C lower than at present during the Last Glacial Maximum. This temperature drop alone would however not have been enough to generate widespread glaciation or permafrost in the Drakensberg Mountains or the Lesotho Highlands. Seasonal freezing of the ground in the Lesotho Highlands might have reached depths of 2 meters or more below the surface. A few small glaciers did however develop during the LGM, in particular in south-facing slopes. In the Hex River Mountains, in the Western Cape, block streams and terraces found near the summit of Matroosberg evidences past periglacial activity which likely occurred during the LGM. Palaeoclimatological proxies indicate the region around Boomplaas Cave was wetter, with increased winter precipitation. The region of the Zambezi River catchment was colder relative to present and the local drop in mean temperature was seasonally uniform.
On the island of Mauritius in the Mascarenhas Archipelago, open wet forest vegetation dominated, contrasting with the dominantly closed-stratified-tall-forest state of Holocene Mauritian forests.
Asia
There were ice sheets in modern Tibet (although scientists continue to debate the extent to which the Tibetan Plateau was covered with ice) as well as in Baltistan and Ladakh. In Southeast Asia, many smaller mountain glaciers formed, and permafrost covered Asia as far south as Beijing. Because of lowered sea levels, many of today's islands were joined to the continents: the Indonesian islands as far east as Borneo and Bali were connected to the Asian continent in a landmass called Sundaland. Palawan was also part of Sundaland, while the rest of the Philippine Islands formed one large island separated from the continent only by the Sibutu Passage and the Mindoro Strait.
The environment along the coast of South China was not very different from that of the present day, featuring moist subtropical evergreen forests, despite sea levels in the South China Sea being about 100 metres lower than the present day.
Australasia
The Australian mainland, New Guinea, Tasmania and many smaller islands comprised a single land mass. This continent is now sometimes referred to as Sahul. In the Bonaparte Gulf of northwestern Australia, sea levels were about 125 metres lower than present. Interior Australia saw widespread aridity, evidenced by extensive dune activity and falling lake levels. Eastern Australia experienced two nadirs in temperature. Lacustrine sediments from North Stradbroke Island in coastal Queensland indicated humid conditions. Data from Little Llangothlin Lagoon likewise indicate the persistence of rainforests in eastern Australia at this time. Rivers maintained their sinuous form in southeastern Australia and there was increased aeolian deposition of sediment in compared to today. The Flinders Ranges likewise experienced humid conditions. In southwestern Western Australia, forests disappeared during the LGM.
Between Sahul and Sundaland – a peninsula of South East Asia that comprised present-day Malaysia and western and northern Indonesia – there remained an archipelago of islands known as Wallacea. The water gaps between these islands, Sahul and Sundaland were considerably narrower and fewer in number than in the present day.
The two main islands of New Zealand, along with associated smaller islands, were joined as one landmass. Virtually all of the Southern Alps were under permanent ice cover, with alpine glaciers extending from them into much of the surrounding high country.
Europe
Northern Europe was largely covered by ice, with the southern boundary of the ice sheets passing through Germany and Poland. This ice extended northward to cover Svalbard and Franz Josef Land and northeastward to occupy the Barents Sea, the Kara Sea, and Novaya Zemlya, ending at the Taymyr Peninsula in what is now northwestern Siberia. Warming commenced in northern latitudes around 20,000 years ago, but it was limited and considerable warming did not take place until around 14,600 year ago.
In northwestern Russia, the Fennoscandian ice sheet reached its LGM extent approximately 17,000 years ago, about five thousand years later than in Denmark, Germany and Western Poland. Outside the Baltic Shield, and in Russia in particular, the LGM ice margin of the Fennoscandian Ice Sheet was highly lobate. The main LGM lobes of Russia followed the Dvina, Vologda and Rybinsk basins respectively. Lobes originated as result of ice following shallow topographic depressions filled with a soft sediment substrate. The northern Ural region was covered in periglacial steppes.
Permafrost covered Europe south of the ice sheet down to as far south as present-day Szeged in Southern Hungary. Ice covered the whole of Iceland. In addition, ice covered Ireland along with roughly the northern half of the British Isles with the southern boundary of the ice sheet running approximately from the south of Wales to the north east of England, and then across the now submerged land of Doggerland to Denmark. Central Europe had isolated pockets of relative warmth corresponding to hydrothermally active areas, which served as refugia for taxa not adapted to extremely cold climates.
In the Cantabrian Mountains of the northwestern corner of the Iberian Peninsula, which in the present day have no permanent glaciers, the LGM led to a local glacial recession as a result of increased aridity caused by the growth of other ice sheets farther to the east and north, which drastically limited annual snowfall over the mountains of northwestern Spain. The Cantabrian alpine glaciers had previously expanded between approximately 60,000 and 40,000 years ago during a local glacial maximum in the region.
In northeastern Italy, in the region around Lake Fimon, Artemisia-dominated semideserts, steppes, and meadow-steppes replaced open boreal forests at the start of the LGM, specifically during Heinrich Stadial 3. The overall climate of the region became both drier and colder.
In the Sar Mountains, the glacial equilibrium-line altitude was about 450 metres lower than in the Holocene. In Greece, steppe vegetation predominated.
Megafaunal abundance in Europe peaked around 27,000 and 21,000 BP; this bountifulness was attributable to the cold stadial climate.
North America
In Greenland, the difference between LGM temperatures and present temperatures was twice as great during winter as during summer. Greenhouse gas and insolation forcings dominated temperature changes in northern Greenland, whereas Atlantic meridional overturning circulation (AMOC) variability was the dominant influence on southern Greenland's climate. Illorsuit Island was exclusively covered by cold-based glaciers.
Eastern Beringia was extremely cold and dry. July air temperatures in northern Alaska and Yukon were about 2-3 °C lower compared to today. Equilibrium line altitudes in Alaska suggest summer temperatures were 2-5 °C compared to preindustrial. Sediment core analysis from Lone Spruce Pond in southwestern Alaska show it was a pocket of relative warmth.
Following a preceding period of relative retreat from 52,000 to 40,000 years ago, the Laurentide Ice Sheet grew rapidly at the onset of the LGM until it covered essentially all of Canada east of the Rocky Mountains and extended roughly to the Missouri and Ohio Rivers, and eastward to Manhattan, reaching a total maximum volume of around 26.5 to 37 million cubic kilometres. At its peak, the Laurentide Ice Sheet reached 3.2 km in height around Keewatin Dome and about 1.7-2.1 km along the Plains divide. In addition to the large Cordilleran Ice Sheet in Canada and Montana, alpine glaciers advanced and (in some locations) ice caps covered much of the Rocky and Sierra Nevada Mountains further south. Latitudinal gradients were so sharp that permafrost did not reach far south of the ice sheets except at high elevations. Glaciers forced the early human populations who had originally migrated from northeast Siberia into refugia, reshaping their genetic variation by mutation and drift. This phenomenon established the older haplogroups found among Native Americans, and later migrations are responsible for northern North American haplogroups.
In southeastern North America, between the southern Appalachian Mountains and the Atlantic Ocean, there was an enclave of unusually warm climate.
South America
In the Southern Hemisphere, the Patagonian Ice Sheet covered the whole southern third of Chile and adjacent areas of Argentina. On the western side of the Andes the ice sheet reached sea level as far north as in the 41 degrees south at Chacao Channel. The western coast of Patagonia was largely glaciated, but some authors have pointed out the possible existence of ice-free refugia for some plant species. On the eastern side of the Andes, glacier lobes occupied the depressions of Seno Skyring, Seno Otway, Inútil Bay, and Beagle Channel. On the Straits of Magellan, ice reached as far as Segunda Angostura.
During the LGM, valley glaciers in the southern Andes (38–43° S) merged and descended from the Andes occupying lacustrine and marine basins where they spread out forming large piedmont glacier lobes. Glaciers extended about 7 km west of the modern Llanquihue Lake, but not more than 2 to 3 km south of it. Nahuel Huapi Lake in Argentina was also glaciated by the same time. Over most of the Chiloé Archipelago, glacier advance peaked 26,000 years ago, forming a long north–south moraine system along the eastern coast of Chiloé Island (41.5–43° S). By that time the glaciation at the latitude of Chiloé was of ice sheet type contrasting to the valley glaciation found further north in Chile.
Despite glacier advances much of the area west of Llanquihue Lake was still ice-free during the Last Glacial Maximum. During the coldest period of the Last Glacial Maximum vegetation at this location was dominated by Alpine herbs in wide open surfaces. The global warming that followed caused a slow change in vegetation towards a sparsely distributed vegetation dominated by Nothofagus species. Within this parkland vegetation Magellanic moorland alternated with Nothofagus forest, and as warming progressed even warm-climate trees began to grow in the area. It is estimated that the tree line was depressed about 1,000 m relative to present day elevations during the coldest period, but it rose gradually until 19,300 years ago. At that time a cold reversal caused a replacement of much of the arboreal vegetation with Magellanic moorland and Alpine species. On Isla Grande de Chiloé, Magellanic moorland and closed-canopy Nothofagus forests were both present during the LGM, but the former disappeared by the late LGM.
Little is known about the extent of glaciers during Last Glacial Maximum north of the Chilean Lake District. To the north, in the dry Andes of Central and the Last Glacial Maximum is associated with increased humidity and the verified advance of at least some mountain glaciers. Montane glaciers in the northern Andes reached their peak extent approximately 27,000 years ago. In northwestern Argentina, pollen deposits record the altitudinal descent of the treeline during the LGM.
Amazonia was much drier than in the present. δD values from plant waxes from the LGM are significantly more enriched than those in the present and those dating back to MIS 3, evidencing this increased aridity. Eastern Brazil was also affected; the site of Guanambi in Bahia was much drier than today.
Atlantic Ocean
AMOC was weaker and more shallow during the LGM. Sea surface temperatures in the western subtropical gyre of the North Atlantic were around 5 °C colder compared to today. Intermediate depth waters of the North Atlantic were better ventilated during the LGM by Glacial North Atlantic Intermediate Water (GNAIW) relative to its present-day ventilation by upper North Atlantic Deep Water (NADW). GNAIW was nutrient poor compared to present day upper NADW. Below GNAIW, southern source bottom water that was very rich in nutrients filled the deep North Atlantic.
Due to the presence of immense ice sheets in Europe and North America, continental weathering flux into the North Atlantic was reduced, as measured by the increased proportion of radiogenic isotopes in neodymium isotope ratios.
There is controversy whether upwelling off the Moroccan coast was stronger during the LGM compared to today. Though coccolith size increases in Calcidiscus leptoporus suggest stronger trade winds during the LGM caused there to be increased coastal upwelling of the northwestern coast of Africa, planktonic foraminiferal δ13C records show upwelling and primary productivity were not enhanced during the LGM except in transient intervals around 23,200 and 22,300 BP.
In the western South Atlantic, where Antarctic Intermediate Water forms, sinking particle flux was heightened as a result of increased dust flux during the LGM and sustained export productivity. The increased sinking particle flux removed neodymium from shallow waters, producing an isotopic ratio change.
Pacific Ocean
On the Island of Hawaii, geologists have long recognized deposits formed by glaciers on Mauna Kea during recent ice ages. The latest work indicates that deposits of three glacial episodes since 150,000 to 200,000 years ago are preserved on the volcano. Glacial moraines on the volcano formed about 70,000 years ago and from about 40,000 to 13,000 years ago. If glacial deposits were formed on Mauna Loa, they have long since been buried by younger lava flows.
Low sea surface temperature (SST) and sea surface salinity (SSS) in the East China Sea during the LGM suggests the Kuroshio Current was reduced in strength relative to the present. Abyssal Pacific overturning was weaker during the LGM than in the present day, although it was temporarily stronger during some intervals of ice sheet retreat. The El Niño–Southern Oscillation (ENSO) was strong during the LGM. Evidence suggests that the Peruvian Oxygen Minimum Zone in the eastern Pacific was weaker than it is in the present day, likely as a result of increased oxygen concentrations in seawater permitted by cooler ocean water temperatures, though it was similar in spatial extent.
The outflow of North Pacific Intermediate Water through the Tasman Sea was stronger during the LGM.
In the Great Barrier Reef along the coast of Queensland, reef development shifted seaward due to the precipitous drop in sea levels, reaching a maximum distance from the present coastline as sea levels approached their lowest levels around 20,700-20,500 years ago. Microbial carbonate deposition in the Great Barrier Reef was enhanced due to low atmospheric CO2 levels.
Indian Ocean
The deep waters of the Indian Ocean were significantly less oxygenated during the LGM compared to the Middle Holocene. The deep South Indian Ocean in particular was an enormous carbon sink, partially explaining the very low pCO2 of the LGM. The intermediate waters of the southeastern Arabian Sea were poorly ventilated relative to today because of the weakened thermohaline circulation.
Southern Ocean
Evidence from sediment cores in the Scotia Sea suggests the Antarctic Circumpolar Current was weaker during the LGM than during the Holocene. The Antarctic Polar Front (APF) was located much farther to the north compared to its present-day location. Studies suggest it could have been placed as far north as 43°S, reaching into the southern Indian Ocean.
Late Glacial Period
The Late Glacial Period followed the LGM and preceded the Holocene, which started around 11,700 years ago.
| Physical sciences | Events | Earth science |
1179005 | https://en.wikipedia.org/wiki/Cerulean | Cerulean | The color cerulean (American English) or caerulean (British English, Commonwealth English), is a variety of the hue of blue that may range from a light azure blue to a more intense sky blue, and may be mixed as well with the hue of green. The first recorded use of cerulean as a color name in English was in 1590. The word is derived from the Latin word caeruleus (), "dark blue, blue, or blue-green", which in turn probably derives from caerulum, diminutive of caelum, "heaven, sky".
"Cerulean blue" is the name of a blue-green pigment consisting of cobalt stannate (). The pigment was first synthesized in the late eighteenth century by Albrecht Höpfner, a Swiss chemist, and it was known as Höpfner blue during the first half of the nineteenth century. Art suppliers began referring to cobalt stannate as cerulean in the second half of the nineteenth century. It was not widely used by artists until the 1870s when it became available in oil paint.
Pigment characteristics
The primary chemical constituent of the pigment is cobalt(II) stannate (). The pigment is a greenish-blue color. In watercolor, it has a slight chalkiness. When used in oil paint, it loses this quality.
Today, cobalt chromate is sometimes marketed under the cerulean blue name but is darker and greener than the cobalt stannate version. The chromate makes excellent turquoise colors and is identified by Rex Art and some other manufacturers as "cobalt turquoise".
Cerulean is inert with good light resistance, and it exhibits a high degree of stability in both watercolor and acrylic paint.
History
Cobalt stannate pigment was first synthesized in 1789 by the Swiss chemist Albrecht Höpfner by heating roasted cobalt and tin oxides together. Subsequently, there was limited German production under the name of Cölinblau. It was generally known as Höpfner blue from the late eighteenth century until the middle of the nineteenth century.
In the late 1850s, art suppliers begin referring to the pigment as "ceruleum" blue. The London Times of 28 December 1859 had an advertisement for "Caeruleum, a new permanent color prepared for the use of artists." Ure's Dictionary of Arts from 1875 describes the pigment as "Caeruleum . . . consisting of stannate of protoxide of cobalt, mixed with stannic acid and sulphate of lime." Cerulean was also referred to as coeurleum, cerulium, bleu céleste (celestial blue). Other nineteenth century English pigment names included "ceruleum blue" and "corruleum blue". By 1935, Max Doerner referred to the pigment as cerulean, as do most modern sources, though ceruleum is still used.
Some sources claim that cerulean blue was first marketed in the United Kingdom by colourman George Rowney, as "coeruleum" in the early 1860s. However, the British firm of Roberson was buying "Blue No. 58 (Cerulium)" from a German firm of Frauenknecht and Stotz prior to Rowney. Cerulean blue was only available as a watercolor in the 1860s and was not widely adopted until the 1870s when it was used in oil paint. It was popular with artists including Claude Monet, Paul Signac, and Picasso. Van Gogh created his own approximation of cerulean blue using a mixture of cobalt blue, cadmium yellow, and white.
Notable occurrences
In 1877, Monet had added the pigment to his palette, using it in a painting from his series La Gare Saint-Lazare (now in the National Gallery, London). The blues in the painting include cobalt and cerulean blue, with some areas of ultramarine. Laboratory analysis conducted by the National Gallery identified a relatively pure example of cerulean blue pigment in the shadows of the station's canopy. Researchers at the National Gallery suggested that "cerulean probably offered a pigment of sufficiently greenish tone to displace Prussian blue, which may not have been popular by this time."
Berthe Morisot painted the blue coat of the woman in her Summer's Day, 1879 in cerulean blue in conjunction with artificial ultramarine and cobalt blue.
When the United Nations was formed at the end of World War II, they adopted cerulean blue for their emblem. The designer Oliver Lundquist stated that he chose the color because it was "the opposite of red, the color of war."
In the Catholic Church, cerulean vestments are permitted on certain Marian feast days, primarily the Immaculate Conception in diocese currently or formerly under the Spanish Crown.
Other color variations
Pale cerulean
Pantone, in a press release, declared the pale hue of cerulean at right, which they call cerulean, as the "color of the millennium".
The source of this color is the "Pantone Textile Paper eXtended (TPX)" color list, color #15-4020 TPX—Cerulean.
Cerulean (Crayola)
This bright tone of cerulean is the color called cerulean by Crayola crayons.
Cerulean frost
At right is displayed the color cerulean frost.
Cerulean frost is one of the colors in the special set of metallic colored Crayola crayons called Silver Swirls, the colors of which were formulated by Crayola in 1990.
Curious Blue
Curious Blue is one of the brighter-toned colors of cerulean.
In nature
Cerulean cuckooshrike
Cerulean kingfisher
Cerulean flycatcher
Cerulean warbler
Cerulean-capped manakin
| Physical sciences | Colors | Physics |
1179089 | https://en.wikipedia.org/wiki/Sarcosuchus | Sarcosuchus | Sarcosuchus (; ) is an extinct genus of crocodyliform and distant relative of living crocodilians that lived during the Early Cretaceous, from the late Hauterivian to the early Albian, 133 to 112 million years ago of what is now Africa and South America. The genus name comes from the Greek σάρξ (sarx) meaning flesh and σοῦχος (souchus) meaning crocodile. It was one of the largest pseudosuchians, with the largest specimen of S. imperator reaching approximately long and weighing up to . It is known from two species; S. imperator from the early Albian Elrhaz Formation of Niger, and S. hartti from the Late Hauterivian of northeastern Brazil. Other material is known from Morocco and Tunisia and possibly Libya and Mali.
The first remains were discovered during several expeditions led by the French paleontologist Albert-Félix de Lapparent, spanning from 1946 to 1959, in the Sahara. These remains were fragments of the skull, vertebrae, teeth, and scutes. In 1964, an almost complete skull was found in Niger by the French CEA, but it was not until 1997 and 2000 that most of its anatomy became known to science, when an expedition led by the American paleontologist Paul Sereno discovered six new specimens, including one with about half the skeleton intact and most of the spine.
Description
Sarcosuchus is a distant relative of living crocodilians, with fully grown individuals estimated to have reached up to in total length and in weight. It had somewhat telescoped eyes and a long snout comprising 75% of the length of the skull. There were 35 teeth in each side of the upper jaw, while in the lower jaw there were 31 teeth in each side. The upper jaw was also noticeably longer than the lower one, leaving a gap between them when the jaws were shut that created an overbite. In young individuals the shape of the snout resembled that of the living gharial, but in fully grown individuals it became considerably broader.
Bulla
Sarcosuchus has an expansion at the end of its snout known as a bulla, which has been compared with the ghara seen in gharials. However, unlike the ghara, which is only found in male gharials, the bulla is present in all Sarcosuchus skulls that have been found so far, suggesting that it was not a sexually dimorphic trait. The purpose of this structure is not known.
Osteoderms
The osteoderms, also known as dermal scutes, of Sarcosuchus were similar to those goniopholodids like Sunosuchus and Goniopholis; they formed an uninterrupted surface that started in the posterior part of the neck down to the middle of the tail as is seen in Araripesuchus and other basal crocodyliforms; this differs from the pattern seen in living crocodiles, which presents discontinuity between the osteoderms of the neck and body.
Size
A common method to estimate the size of crocodiles and crocodile-like reptiles is the use of the length of the skull measured in the midline from the tip of the snout to the back of the skull table, as in living crocodilians there is a strong correlation between skull length and total body length in subadult and adult individuals irrespective of their sex. This method was used by Sereno et al. (2001) for Sarcosuchus due to the absence of a complete enough skeleton. Two regression equations were used to estimate the size of S. imperator, they were created based on measurements gathered from 17 captive gharial individuals from northern India and from 28 wild saltwater crocodile individuals from northern Australia, both datasets supplemented by available measurements of individuals over in length found in the literature. The largest known skull of S. imperator (the type specimen) is long ( in the midline), and it was estimated that the individual it belonged to had a total body length of . Its snout-vent length of was estimated using linear equations for the saltwater crocodile and in turn this measurement was used to estimate its body weight at . This shows that Sarcosuchus was able to reach a maximum body size not only greater than previously estimated but also greater than that of the Miocene Rhamphosuchus, the Late Cretaceous Deinosuchus, and the Miocene Purussaurus according to current estimates at that time.
However, extrapolation from the femur of a subadult individual as well as measurements of the skull width further showed that the largest S. imperator was significantly smaller than was estimated by Sereno et al. (2001) based on modern crocodilians. O’Brien et al. (2019) estimated the length of the largest S. imperator specimen at nearly and body mass at based on longirostrine crocodylian skull width to total length and body width ratio. The highest upper quartile reconstructed length and body mass for the specimen is and , respectively.
Classification
Sarcosuchus is commonly classified as part of the clade Pholidosauridae, a group of crocodile-like reptiles (Crocodyliformes) related but outside Crocodylia (the clade containing living crocodiles, alligators and gharials). Within this group it is most closely related to the North American genus Terminonaris. Most members of Pholidosauridae had long, slender snouts and they all were aquatic, inhabiting several different environments. Some forms are interpreted as marine, capable of tolerating saltwater while others, like Sarcosuchus, were freshwater forms. The most primitive members of the clade, however, were found in coastal settings, zones mixing freshwater and marine waters. Sarcosuchus stands out among pholidosaurids for being considered a generalist predator, different from most known members of the clade which were specialized piscivores. A 2019 study found it to be in a more derived position in Tethysuchia, being phylogenetically closer to Dyrosauridae.
Simplified cladogram after Fortier et al. (2011).
Discovery and naming
Early findings
During the course of several expeditions on the Sahara from 1946 to 1959 which were led by the French paleontologist Albert-Félix de Lapparent, several fossils of a crocodyliform of large size were unearthed in the region known as the Continental Intercalaire Formation. Some of them were found in Foggara Ben Draou, in Mali and near the town of Aoulef, Algeria (informally named as the Aoulef Crocodile) while others came from the Ain el Guettar Formation of Gara Kamboute. In the south of Tunisia, the fossils found were fragments of the skull, teeth, scutes and vertebrae. In 1957, in the region now known as the Elrhaz Formation, several isolated teeth of great size were found by H. Faure. The study of this material by French paleontologist France De Broin helped identify them as coming from a long-snouted crocodile.
Later, in 1964, the research team of the French CEA discovered an almost complete skull in the region of Gadoufaoua in the Niger. The said skull was shipped to Paris for study and became the holotype of the then new genus and species Sarcosuchus imperator in 1966.
Fossils from Brazil
In 1977, a new species of Sarcosuchus was recognised, S. hartti, from remains found in the late 19th century in late Hauterivian pebbly conglomerates and green shales belonging to the Ilhas Formation in the Recôncavo Basin of north-eastern Brazil. In 1867, American naturalist Charles Hartt found two isolated teeth and sent them to the American paleontologist O. C. Marsh who erected a new species of Crocodylus for them, C. hartti. This material, along with other remains were assigned in 1907 to the genus Goniopholis as G. hartti. Now residing in the British Museum of Natural History, the fragment of the lower jaw, dorsal scute and two teeth compromising the species G. hartti were reexamined and conclusively placed in the genus Sarcosuchus.
Recent findings
The next major findings occurred during the expeditions led by the American paleontologist Paul Sereno in 1995 (Aoufous Formation, Morocco), 1997 and the follow-up trip in 2000. Partial skeletons, numerous skulls and 20 tons of assorted other fossils were recovered from the deposits of the Elrhaz Formation, which has been dated as late Aptian or early Albian stages of the Late Cretaceous. It took about a year to prepare the Sarcosuchus remains.
A tooth enamel from the Ifezouane Formation (lower Kem Kem beds) of Morocco was identified as cf. Sarcosuchus. Fossil teeth from the area of Nalut in northwestern Libya, possibly Hauterivian to Barremian in age, might be referable to S. imperator. Indeterminate Sarcosuchus material including dorsal osteoderms in anatomical connection, isolated teeth and fragmentary skeletal remains including a left scapula, mandible fragment, dorsal vertebrae, ilium and a proximal portion of a femur was described from the Oum Ed Dhiab Member in Tunisia in 2018.
Paleobiology
Growth pattern
Sereno took thin sections from trunk osteoderms of an estimated subadult individual (~80% of estimated maximum adult size). Approximately 40 lines of arrested growth (LAG) were counted in these thin sections, suggesting that S. imperator took 50 to 60 years to reach adult size. Given that extant wild crocodylians rarely reach these advanced ages, Sereno suggested that S. imperator achieved its large size by extending its period of rapid, juvenile, growth. A similar growth strategy has been suggested for the equally titanic crocodylian Deinosuchus, based on similar criteria.
Diet
Based on the broader snout of fully grown S. imperator when compared with the living gharial and other narrow-snouted crocodiles, along with a lack of interlocking of the smooth and sturdy-crowned teeth when the jaws were closed, Sereno et al. hypothesized that S. imperator had a generalized diet similar to that of the Nile crocodile, which would have included large terrestrial prey such as the abundant dinosaurs that lived in the same region.
However, a 2014 analysis of a biomechanical model of its skull suggested that unlike Deinosuchus, Sarcosuchus may not have been able to perform the "death roll" maneuver used by extant crocodilians to dismember their prey. This suggests that if S. imperator did hunt big game, it probably did not dismember prey in the same fashion as extant crocodilians.
Habitat
The remains of S. imperator were found in a region of the Ténéré Desert named Gadoufaoua, more specifically in the Elrhaz Formation of the Tegama Group, dating from the late Aptian to the early Albian of the Early Cretaceous, approximately 112 million years ago. The stratigraphy of the region and the aquatic fauna that was found therein indicates that it was an inland fluvial environment, entirely freshwater in nature with a humid tropical climate. S. imperator shared the waters with the holostean fish Lepidotus and the coelacanth Mawsonia. The dinosaur fauna was represented by the iguanodontian Lurdusaurus, which was the most common dinosaur in the region, and its relative Ouranosaurus; there were also two sauropods, Nigersaurus and a currently unnamed sauropod while the theropod fauna included the spinosaurid Suchomimus, the carcharodontosaurid Eocarcharia and the abelisaurid Kryptops.
Meanwhile, S. hartti was found in the Recôncavo Basin of Brazil, specifically in the Ilhas Formation of the Bahia series. It was a shallow lacustrine environment dating from the late Aptian, similar in age to the habitat of S. imperator, with similar aquatic fauna, including Lepidotus and two species of Mawsonia. The dinosaur fauna is of a very fragmentary nature and identification does not go beyond indeterminate theropod and iguanodontid remains.
| Biology and health sciences | Prehistoric crocodiles | Animals |
1179384 | https://en.wikipedia.org/wiki/Java%20sparrow | Java sparrow | The Java sparrow (Lonchura oryzivora; Japanese: 文鳥, bunchō), also known as the Java finch, Java rice sparrow or Java rice bird, is a small passerine bird. This estrildid finch is a resident breeding bird in Java, Bali and Bawean in Indonesia. It is a popular cage bird, and has been introduced into many other countries.
Taxonomy
The Java sparrow was formally described by the Swedish naturalist Carl Linnaeus in 1758 in the tenth edition of his Systema Naturae under the binomial name Loxia oryzivora. The specific epithet combines Latin oryza meaning "rice" with -vorus meaning "eating". Linnaeus based his description on the "Padda or Rice-bird" that had been described and illustrated in 1743 by the English naturalist George Edwards in his A Natural History of Uncommon Birds. Edwards believed that his specimens had come from China but mentions the common name "Java sparrow". The species was reclassified to the genus Lonchura in 2020.
Description
The Java sparrow is about in length from the beak to its tip of tail feathers. Although only about the size of a house sparrow, it may be the largest species in the estrildid family. The mean body mass is , making it slightly heavier than its nearest known rival, the black-bellied seedeater. The adult is unmistakable, with its grey upperparts and breast, pink belly, white-cheeked black head, red eye-ring, pink feet and thick red bill.
Both sexes are similar. Immature birds have brown upperparts and pale brown underparts, and a plain head. Very young birds have a black beak with a pink base.
The call is a chip, and the song is a rapid series of call notes chipchipchipchipchipchip.
Java sparrows produce distinct trill-calls in different behavioral contexts, according to a study by Furutani et al. (2018). These trill-calls, though acoustically similar, vary based on their repetition rate, which changes depending on the situation. In aggressive interactions, the sparrows emit faster trill-calls with higher sound pressure levels and entropy. In contrast, during affiliative behaviors, the trill-calls are slower and softer. This variation in trill-calls plays a crucial role in the birds' social communication, helping them convey different intentions based on the context.
Habitat
The Java sparrow is a very gregarious bird which feeds mainly on grain and other seeds. It frequents open grassland and cultivation, and was formerly a pest in rice fields, hence its scientific name. The nest is constructed in a tree or building, and up to eight eggs are laid.
Aviculture
The Java sparrow has been a popular cage bird in Asia for centuries, first in China's Ming Dynasty and then in Japan from the 17th century, frequently appearing in Japanese paintings and prints. Meiji-era writer Natsume Sōseki wrote an essay about his pet Java sparrow. In the late 1960s and early 1970s the Java sparrow was one of the most popular cage birds in the United States until its import was banned. Today it remains illegal to possess in California because of a perceived threat to agriculture, although rice-dependent Asian countries like China, Taiwan and Japan have not regulated the bird.
In Asia the Java sparrow is most often raised almost from birth by human breeders and owners, and they become very tame and attached to humans. As such, they can be normally kept in relatively small cages, but let out for indoor exercise without their attempting to escape. In captivity, a variety of colourations have been bred, including white, silver/opal, fawn/isabel, pastel, cream and agate (which currently is rare within Europe captive specimens) along with the pied Java sparrow (called the sakura buncho in Japan).
Introductions
The Java sparrow was introduced in the Indian subcontinent, but it failed to become a successful resident on the Indian mainland. In the United States there are breeding populations on several of the Hawaiian Islands, especially Oahu.
In the Caribbean, the Java sparrow was introduced to Puerto Rico where it is fairly common near San Juan. It has also been sighted in Jamaica, but is not known to occur on any of the other islands. It was also introduced to Christmas Island, off the coast of Western Australia.
Threats
The Java sparrow is considered by some countries to be an agricultural pest with respect to rice cultivation. An ongoing loss of natural habitat, hunting in some areas and trapping (as a pest) in others has led to much smaller numbers in the wild and sightings in its natural range have become increasingly uncommon. The Java sparrow is now evaluated as endangered on the IUCN Red List of Threatened Species with less than 10,000 individuals remaining. It is also listed on Appendix II of CITES. The species is also severely threatened by the illegal exotic pet trade as they are sought after for their distinctive song, according to TRAFFIC.
| Biology and health sciences | Passerida | Animals |
1179801 | https://en.wikipedia.org/wiki/Boracite | Boracite | Boracite is a magnesium borate mineral with formula: Mg3B7O13Cl. It occurs as blue green, colorless, gray, yellow to white crystals in the orthorhombic - pyramidal crystal system. Boracite also shows pseudo-isometric cubical and octahedral forms. These are thought to be the result of transition from an unstable high temperature isometric form on cooling. Penetration twins are not unusual. It occurs as well formed crystals and dispersed grains often embedded within gypsum and anhydrite crystals. It has a Mohs hardness of 7 to 7.5 and a specific gravity of 2.9. Refractive index values are nα = 1.658 - 1.662, nβ = 1.662 - 1.667 and nγ = 1.668 - 1.673. It has a conchoidal fracture and does not show cleavage. It is insoluble in water (not to be confused with borax, which is soluble in water).
Boracite is typically found in evaporite sequences associated with gypsum, anhydrite, halite, sylvite, carnallite, kainite and hilgardite. It was first described in 1789 for specimens from its type locality of Kalkberg hill, Lüneburg, Lower Saxony, Germany. It is also found near Sussex, New Brunswick.
The name is derived from its boron content (19 to 20% boron by mass).
| Physical sciences | Minerals | Earth science |
1179836 | https://en.wikipedia.org/wiki/Enstatite | Enstatite | Enstatite is a mineral; the magnesium endmember of the pyroxene silicate mineral series enstatite (MgSiO3) – ferrosilite (FeSiO3). The magnesium rich members of the solid solution series are common rock-forming minerals found in igneous and metamorphic rocks. The intermediate composition, , has historically been known as hypersthene, although this name has been formally abandoned and replaced by orthopyroxene. When determined petrographically or chemically the composition is given as relative proportions of enstatite (En) and ferrosilite (Fs) (e.g., En80Fs20).
Polymorphs and varieties
Most natural crystals are orthorhombic (space group Pbca) although three polymorphs are known. The high temperature, low pressure polymorphs are protoenstatite and protoferrosilite (also orthorhombic, space group Pbcn) while the low temperature forms, clinoenstatite and clinoferrosilite, are monoclinic (space group P21/c).
Weathered enstatite with a small amount of iron takes on a submetallic luster and a bronze-like color. This material is termed bronzite, although it is more correctly called altered enstatite.
Bronzite and hypersthene were known long before enstatite, which was first described by G. A. Kenngott in 1855.
An emerald-green variety of enstatite is called chrome-enstatite and is cut as a gemstone. The green color is caused by traces of chromium, hence the varietal name. In addition, black, chatoyant hypersthene and brownish bronzite are also used as semi-precious gemstones.
Identification
Enstatite and the other orthorhombic pyroxenes are distinguished from those of the monoclinic series by their optical characteristics, such as straight extinction, much weaker double refraction and stronger pleochroism. They also have a prismatic cleavage that is perfect in two directions at 90 degrees. Enstatite is white, gray, greenish, or brown in color; its hardness is 5–6 on the Mohs scale, and its specific gravity is 3.2–3.3. This prismatic form is used in gemstones, and for academic purposes.
Occurrence
Isolated crystals are rare, but orthopyroxene is an essential constituent of various types of igneous rocks and metamorphic rocks. Magnesian orthopyroxene occurs in plutonic rocks such as gabbro (norite) and diorite. It may form small idiomorphic phenocrysts and also groundmass grains in volcanic rocks such as basalt, andesite, and dacite.
Enstatite, close to En90Fs10 in composition, is an essential mineral in typical peridotite and pyroxenite of the Earth's mantle. Xenoliths of peridotite are common in kimberlite and in some basalt. Measurements of the calcium, aluminum, and chromium contents of enstatite in these xenoliths have been crucial in reconstructing the depths from which the xenoliths were plucked by the ascending magmas.
Orthopyroxene is an important constituent of some metamorphic rocks such as granulite. Orthopyroxene near pure enstatite in composition occurs in some metamorphosed serpentines. Large crystals, a foot in length and mostly altered to steatite, were found in 1874 in the apatite veins traversing mica-schist and hornblende-schist at the apatite mine of Kjørstad, near Brevik in southern Norway.
Enstatite is a common mineral in meteorites. Crystals have been found in stony and iron meteorites, including one that fell at Breitenbach in the Ore Mountains, Bohemia. In some meteorites, together with olivine it forms the bulk of the material; it can occur in small spherical masses, or chondrules, with an internal radiated structure.
In space
Enstatite is one of the few silicate minerals that have been observed in crystalline form outside the Solar System, particularly around evolved stars and planetary nebulae such as NGC 6302. Enstatite is thought to be one of the early stages for the formation of crystalline silicates in space. Many correlations have been noted between the occurrence of the mineral and the structure of the object around which it has been observed.
Enstatite is thought to be a main component of the E-type asteroids. The Hungaria asteroids are the main examples in the Solar System.
A layer of quartz and enstatite clouds above an iron cloud deck are thought to exist in the atmosphere of the young brown dwarf 2M2224-0158.
| Physical sciences | Silicate minerals | Earth science |
228385 | https://en.wikipedia.org/wiki/Clock%20rate | Clock rate | Clock rate or clock speed in computing typically refers to the frequency at which the clock generator of a processor can generate pulses used to synchronize the operations of its components. It is used as an indicator of the processor's speed. Clock rate is measured in the SI unit of frequency hertz (Hz).
The clock rate of the first generation of computers was measured in hertz or kilohertz (kHz), the first personal computers from the 1970s through the 1980s had clock rates measured in megahertz (MHz). In the 21st century the speed of modern CPUs is commonly advertised in gigahertz (GHz). This metric is most useful when comparing processors within the same family, holding constant other features that may affect performance.
Determining factors
Binning
Manufacturers of modern processors typically charge higher prices for processors that operate at higher clock rates, a practice called binning. For a given CPU, the clock rates are determined at the end of the manufacturing process through testing of each processor. Chip manufacturers publish a "maximum clock rate" specification, and they test chips before selling them to make sure they meet that specification, even when executing the most complicated instructions with the data patterns that take the longest to settle (testing at the temperature and voltage that gives the lowest performance). Processors successfully tested for compliance with a given set of standards may be labeled with a higher clock rate, e.g., 3.50 GHz, while those that fail the standards of the higher clock rate yet pass the standards of a lower clock rate may be labeled with the lower clock rate, e.g., 3.3 GHz, and sold at a lower price.
Engineering
The clock rate of a CPU is normally determined by the frequency of an oscillator crystal. Typically a crystal oscillator produces a fixed sine wave—the frequency reference signal. Electronic circuitry translates that into a square wave at the same frequency for digital electronics applications (or, when using a CPU multiplier, some fixed multiple of the crystal reference frequency). The clock distribution network inside the CPU carries that clock signal to all the parts that need it. An A/D Converter has a "clock" pin driven by a similar system to set the sampling rate. With any particular CPU, replacing the crystal with another crystal that oscillates at half the frequency ("underclocking") will generally make the CPU run at half the performance and reduce waste heat produced by the CPU. Conversely, some people try to increase performance of a CPU by replacing the oscillator crystal with a higher frequency crystal ("overclocking"). However, the amount of overclocking is limited by the time for the CPU to settle after each pulse, and by the extra heat created.
After each clock pulse, the signal lines inside the CPU need time to settle to their new state. That is, every signal line must finish transitioning from 0 to 1, or from 1 to 0. If the next clock pulse comes before that, the results will be incorrect. In the process of transitioning, some energy is wasted as heat (mostly inside the driving transistors). When executing complicated instructions that cause many transitions, the higher the clock rate the more heat produced. Transistors may be damaged by excessive heat.
There is also a lower limit of the clock rate, unless a fully static core is used.
Historical milestones and current records
The first fully mechanical digital computer, the Z1, operated at 1 Hz (cycle per second) clock frequency and the first electromechanical general purpose computer, the Z3, operated at a frequency of about 5–10 Hz. The first electronic general purpose computer, the ENIAC, used a 100 kHz clock in its cycling unit. As each instruction took 20 cycles, it had an instruction rate of 5 kHz.
The first commercial PC, the Altair 8800 (by MITS), used an Intel 8080 CPU with a clock rate of 2 MHz (2 million cycles per second). The original IBM PC (c. 1981) had a clock rate of 4.77 MHz (4,772,727 cycles per second). In 1992, both Hewlett-Packard and Digital Equipment Corporation (DEC) exceeded 100 MHz with RISC techniques in the PA-7100 and AXP 21064 DEC Alpha respectively. In 1995, Intel's P5 Pentium chip ran at 100 MHz (100 million cycles per second). On March 6, 2000, AMD demonstrated passing the 1 GHz milestone a few days ahead of Intel shipping 1 GHz in systems. In 2002, an Intel Pentium 4 model was introduced as the first CPU with a clock rate of 3 GHz (three billion cycles per second corresponding to ~ 0.33 nanoseconds per cycle). Since then, the clock rate of production processors has increased more slowly, with performance improvements coming from other design changes.
Set in 2011, the Guinness World Record for the highest CPU clock rate is 8.42938 GHz with an overclocked AMD FX-8150 Bulldozer-based chip in an LHe/LN2 cryobath, 5 GHz on air. This is surpassed by the CPU-Z overclocking record for the highest CPU clock rate at 8.79433 GHz with an AMD FX-8350 Piledriver-based chip bathed in LN2, achieved in November 2012. It is also surpassed by the slightly slower AMD FX-8370 overclocked to 8.72 GHz which tops off the HWBOT frequency rankings. These records were broken in late 2022 when an Intel Core i9-13900K was overclocked to 9.008 GHz.
The highest base clock rate on a production processor is the i9-14900KS, clocked at 6.2 GHz, which was released in Q1 2024.
Research
Engineers continue to find new ways to design CPUs that settle a little more quickly or use slightly less energy per transition, pushing back those limits, producing new CPUs that can run at slightly higher clock rates. The ultimate limits to energy per transition are explored in reversible computing.
The first fully reversible CPU, the Pendulum, was implemented using standard CMOS transistors in the late 1990s at the Massachusetts Institute of Technology.
Engineers also continue to find new ways to design CPUs so that they complete more instructions per clock cycle, thus achieving a lower CPI (cycles or clock cycles per instruction) count, although they may run at the same or a lower clock rate as older CPUs. This is achieved through architectural techniques such as instruction pipelining and out-of-order execution which attempts to exploit instruction level parallelism in the code.
Comparing
The clock rate of a CPU is most useful for providing comparisons between CPUs in the same family. The clock rate is only one of several factors that can influence performance when comparing processors in different families. For example, an IBM PC with an Intel 80486 CPU running at 50 MHz will be about twice as fast (internally only) as one with the same CPU and memory running at 25 MHz, while the same will not be true for MIPS R4000 running at the same clock rate as the two are different processors that implement different architectures and microarchitectures. Further, a "cumulative clock rate" measure is sometimes assumed by taking the total cores and multiplying by the total clock rate (e.g. a dual-core 2.8 GHz processor running at a cumulative 5.6 GHz). There are many other factors to consider when comparing the performance of CPUs, like the width of the CPU's data bus, the latency of the memory, and the cache architecture.
The clock rate alone is generally considered to be an inaccurate measure of performance when comparing different CPUs families. Software benchmarks are more useful. Clock rates can sometimes be misleading since the amount of work different CPUs can do in one cycle varies. For example, superscalar processors can execute more than one instruction per cycle (on average), yet it is not uncommon for them to do "less" in a clock cycle. In addition, subscalar CPUs or use of parallelism can also affect the performance of the computer regardless of clock rate.
| Technology | Computer hardware | null |
228458 | https://en.wikipedia.org/wiki/Muscle%20car | Muscle car | A muscle car is an American-made two-door sports coupe with a powerful engine, generally designed for high-performance driving.
In 1949, General Motors introduced its 88 with the company's OHV Rocket V8 engine, which was previously available only in its luxury Oldsmobile 98. This formula of putting a maker's largest, most powerful engine in a smaller, lighter, more affordable vehicle evolved into the "muscle car" category. Chrysler and Ford quickly followed suit with the Chrysler Saratoga and the Lincoln Capri.
The term "muscle car", which appeared in the mid-1960s, was originally applied to "performance"-oriented street cars produced to fill a newly recognized niche; it entered the general vocabulary through car magazines and automobile marketing and advertising. By the early 1970s, muscle cars included special editions of mass-production cars designed for street and track drag racing. The concept of high performance at lower prices was exemplified by the 1968 Plymouth Road Runner and companion Dodge Super Bee, whose powerful engines drove relatively basic-trimmed intermediate-sized cars that were meant to undercut more expensive, more stylish, and better-appointed models from General Motors and Ford that had come to define the market, such as the Pontiac GTO (1964), 396 Chevrolet Chevelle (1965), 400 Buick Gran Sport (1965), 400 Oldsmobile 442 (1965), as well as the 427 Mercury Comet Cyclone (1964) and 390 Mercury Cyclone (1966).
By some definitions – including those used by Car and Driver, CNBC, Road & Track, and Motor Trend—pony cars such as the Ford Mustang, Chevrolet Camaro, Plymouth Barracuda, Pontiac Firebird, AMC Javelin, and their luxury companions in that large, influential, and lucrative 1960s–70s niche, the Mercury Cougar and Dodge Challenger, could also qualify as "muscle cars" if outfitted with suitable high-performance equipment.
Terminology
Definition
The definition of a muscle car is subjective and endlessly debated, resulting in the term having few universally agreed characteristics:
A large high-performance V8 engine, often in the most powerful configuration offered for a particular model
Rear-wheel drive
Being manufactured in the United States in the 1960s or early 1970s (the specific year range of 1964–1973 is sometimes used)
A relatively lightweight two-door body (though opinions vary as to whether high-performance full-size cars, compacts, and pony cars qualify as muscle cars, and why a two-seat AMC AMX could be, but a two-seat Chevrolet Corvette was not. While some feel that only mid-size cars can be considered muscle cars, this view is not held by the top, industry-defining, enthusiast publications, including Car and Driver, Road & Track, and Motor Trend.
High-powered pony cars are sometimes considered muscle cars, as by the above-mentioned publications, with some exceptional personal luxury cars also regarded by some as qualifying on their merits. In the opposite direction, by the late 1960s a wave of inexpensive, straight-line speed oriented stripped down intermediate sedans offered at prices under as expanded the original definition from a "muscle car" as one offering both performance and some measure of style, accessories, and cachet, and doubled it back toward the drag strip focus of such exceptional early proto-muscle cars as the limited production, factory experimental 1964 Ford Fairlane Thunderbolt.
Sports cars – including those which meet all the above most basic criterion, such as the 1969 ZL-1 Corvette, with an all-aluminum V8 listed at but reported to produce , that slung the car through the traps in 10.89 seconds – are considered muscle cars by some, and not by others. Drag strip-oriented fans see muscle cars as an extension of the hot rodding philosophy of taking a small car and putting a large-displacement engine in it to maximize straight-line speed. However, widespread public acceptance and use of the term, including that exemplified by the Car and Driver, CNBC, Road & Track, and Motor Trend top muscle car lists below, affirm a much broader interpretation as the norm.
"Supercar"
Muscle cars were initially referred to as "supercars" in the United States, such as the 1957 Rambler Rebel, which was described as a "potent mill turned the lightweight Rambler into a veritable supercar." From the mid-1960s to the mid-1970s, "dragstrip bred" mid-size cars equipped with large V8 engines and rear-wheel drive were also referred to as supercars, more often than muscle cars.
In 1966, the supercar became an "industry trend". This was when the four domestic automakers "needed to cash in on the supercar market" with eye-catching, heart-stopping cars. An example of the use of the supercar description for early muscle car models includes the May 1965 Car Life road test of the Pontiac GTO, followed in 1968 with a Car and Driver review of the 1969 American Motors SC/Rambler describing it as ready to compete in "the Supercar street racer gang" market segment, with the initials "SC" signifying SuperCar, and a 1969 Car Life review that included how "Hurst puts American Motors into the Supercar club with the 390 Rogue".
The supercar market segment in the U.S. at the time included unique versions of regular production models that were positioned in several sizes and market segments (such as the "economy supercar"), as well as limited edition, documented dealer-converted vehicles. However, over time the term came to be applied to much, much more expensive and exotic cars, which claimed the name supercar.
History
1950s: Origins
Opinions on the origin of the muscle car vary, but the 1949 Oldsmobile Rocket 88 is cited as the first full-sized muscle car. The Rocket 88 was the first time a powerful V8 engine was available in a smaller and lighter body style (in this case the engine from the larger Oldsmobile 98 with the body from the six-cylinder Oldsmobile 76). The Rocket 88 produced at 3600 rpm and at 1800 rpm and won eight out of ten races in the 1950 NASCAR season. The Rocket 88's Oldsmobile 303 V8 engine, along with the Cadillac 331 engine, also introduced in 1949, are stated to have "launched the modern era of the high-performance V-8."
In 1955, the large-sized Chrysler C-300 - the first in a long, 15-year series of large, expensive, performance-first Chryslers - was introduced that produced from its V8 engine, and it was advertised as "America's Most Powerful Car". Capable of accelerating from 0 to in 9.8 seconds and reaching , the 1955 Chrysler 300 is also recognized as one of the best-handling cars of its era.
The compact-sized 1956 Studebaker Golden Hawk was powered by a Packard V8, the second most powerful engine to the Chrysler 300.
The Rambler Rebel, introduced by American Motors Corporation (AMC) in 1957, is the first mid-sized car to be available with a big-block V8 engine. The Rebel followed most of the muscle car formula including "make 'em go fast as well as cheaply." It is therefore considered by some to be the first muscle car. With a V8 engine producing , its 0–60 mph acceleration of 7.5 seconds made it the fastest stock American sedan at the time. Only the fuel-injected Chevrolet Corvette beat it by half a second.
Early 1960s: Drag racing influences
The popularity and performance of muscle cars grew in the early 1960s. This was when Mopar (Dodge, Plymouth, and Chrysler) and Ford battled for supremacy in drag racing. The 1961 Chevrolet Impala offered an SS package for $53.80, which consisted of a V8 engine producing along with upgraded brakes, tires, and suspension. The 1962 Dodge Dart 413 (nicknamed Max Wedge) had a V8 which produced and could cover the quarter-mile in under 13 seconds.
In 1963, two hundred Ford Galaxie "R-code" cars were factory-built specifically for drag racing, resulting in a full-size car that could cover the quarter-mile in a little over 12 seconds. Upgrades included fiberglass panels, aluminum bumpers, traction bars, and a Ford FE-based racing engine conservatively rated at . The road-legal version of the Galaxie 427 used the "Q-code" engine which produced . The following year, Ford installed the proven 427 "top-oiler" engine in the smaller and lighter Fairlane body, creating the Ford Thunderbolt. The Thunderbolt included several weight-saving measures (including acrylic windows and fibreglass/aluminium body panels and bumpers) and a stock Thunderbolt could cover the quarter-mile in 11.76 seconds. The Thunderbolt was technically road-legal, however, it was considered unsuitable even "for driving to and from the (drag)strip, let alone on the street in everyday use". A total of 111 Thunderbolts were built.
The General Motors competitor to the Thunderbolt was the Z-11 option package for the full-size Chevrolet Impala coupe, of which 57 examples were produced in 1963 only. The Z-11 Impala was powered by a version of the W-series big-block engine, which was officially rated at . With a compression ratio of 13.5:1, the engine required high-octane fuel. The RPOZ-11 package also included weight reduction measures such as an aluminum hood and fenders, the removal of sound-deadening material as well as the deletion of the heater and radio.
In 1964, a drag racing version of the Dodge 330 was created, called the "330 Lightweight". It was powered by a version of the Hemi racing engine which was official rated at , but rumored to have an actual power output higher than this. Weight reduction measures included an aluminium hood as well as lightweight front bumpers, fenders and doors, polycarbonate side windows, and no sound deadening. Like other lightweights of the era, it came with a factory disclaimer: "Designed for supervised acceleration trials. Not recommended for general everyday driving because of the compromises in the all-round characteristics which must be made for this type of vehicle."
Also using the 426 Hemi racing engine was the limited production 1965 Plymouth Satellite 426 Hemi. In 1966, the racing version of the 426 Hemi was replaced by a detuned "Street Hemi" version, also with a size of 426 cu in and an official power rating of ). The 1966 Plymouth Satellite 426 Hemi could run a 13.8-second quarter-mile at and had a base price of $3,850.
1964–1970: Peak muscle car era
Although pure muscle cars often sold in relatively small volumes, manufacturers valued the halo effect of the publicity created by these models. Competition between manufacturers led to a horsepower war that peaked in 1970, with models such as the LS-6 Chevelle advertising as much as .
The Pontiac GTO, a car that captured the public mind and strongly influenced the muscle car era, was introduced in 1964 as an optional package for the intermediate-size Pontiac Tempest. The GTO was developed by Pontiac division president John DeLorean and was initially powered by a V8 engine producing . The success of the GTO led other GM divisions to develop muscle cars based on intermediate-sized platforms: the 1964 Oldsmobile 442, 1964 Chevrolet Chevelle SS, and 1965 Buick Gran Sport.
The AMC V8 engine was enlarged to in 1968, which produced and was first used in the 1968 AMC Rebel SST, AMC Javelin Go-package, and AMC AMX. AMC was a car manufacturing company that made these two incredible cars. AMC only made small economy cars until they hired Dick Teague as a designer, who later became the vice president. AMC was not known as a muscle car maker, but the new AMX and AMC Javelin became top competitors.
As the 1960s progressed, optional equipment and luxury appointments increased in many popular models of "performance-oriented" cars. With the added weight and power-consuming accessories and features, engines had to be more powerful to maintain performance levels, and the cars became more expensive. In response, some "budget" muscle cars began to appear, such as the 1967 Plymouth GTX, the 1968 Plymouth Road Runner, and the 1968 Dodge Super Bee. In 1969, the Plymouth Road Runner was awarded Motor Trend magazine's Car of the Year. With optional performance parts such as intake and exhaust manifolds, upgraded carburetor, and drag-racing tires, the Road Runner had a quarter-mile time of 14.7 seconds at . In this customized form, the cost of the Road Runner was US$3,893.
The Plymouth Barracuda was a pony car that could be turned into a muscle car with the addition of the famed Chrysler 426 Hemi, available as an option beginning in 1968, after debuting in street form two years earlier in the Plymouth Belvedere, Dodge Coronet, and Dodge Charger. Originally based on the smaller compact car body and chassis of the Plymouth Valiant, the Barracuda was also available with a V8 engine producing . It could run a quarter-mile in 13.33 seconds at on the drag strip. The base price was $2,796.00; the price as tested by Hot Rod was $3,652. The related 1970 Plymouth Duster was powered by a V8 engine producing . Performance figures were 0 to in 6.0 seconds and the quarter-mile time of in 14.7 seconds at .
The Chevrolet L72 big-block engine became available in the mid-sized Chevrolet Chevelle in 1969 as the COPO 427 option. The 427 Chevelle could run a 13.3 sec. quarter-mile at . Chevrolet rated the engine at , but the NHRA claimed power output to be . The following year, the "Chevelle SS 454" model was introduced, which used the Chevrolet LS6 big-block engine rated at , the highest factory rating at that time.
The fastest muscle car produced by American Motors was the mid-sized 1970 AMC Rebel "The Machine", which was powered by a engine producing . The Rebel had a 0– time of 6.8 seconds and a quarter-mile run in 14.4 seconds at .
1970s: Decline of the segment
The popularity of muscle cars declined through the early 1970s, due to factors including the Clean Air Act, the fuel crisis, and the increasing cost of insurance for these types of cars. The 1973 oil crisis resulted in rationing of fuel and higher prices. Muscle cars quickly became unaffordable and impractical for many people. In addition, the automobile insurance industry levied surcharges on all high-powered models.
Before the Clean Air Act of 1970, a majority of muscle cars came optioned with high-compression engines (some engines were as high as 11:1), which required high-octane fuel. Prior to 1970, 100-octane fuel was common. However, following the passage of the Clean Air Act of 1970, octane ratings were lowered to 91 (due in part to the removal of lead). Manufacturers reduced the compression ratio of engines, resulting in reduced performance. Simultaneously, efforts to combat air pollution focused Detroit's attention on emissions control rather than increased power outputs.
1980s–1990s: Performance revival
Muscle car performance began a resurgence in the early 1980s with high-output V8 engines introduced for the Ford Mustang GT, Chevrolet Camaro Z28, and Pontiac Firebird Formula/Trans Am. Initially using four-barrel carburetors, engine performance, and fuel economy were increased by the mid-1980s using electronic fuel injection systems and advanced engine management controls. Muscle car performance began to reappear on intermediate two-door coupés such as the Chevrolet Monte Carlo SS and Buick Regal. The Buick Regal used turbocharged V6 engines on the Grand National, Turbo-T, T-Type, and GNX models which rivaled the performance of V8 engines.
The few muscle cars remaining in production by the mid-1990s included the fourth generation Mustang, fourth generation Camaro, and fourth generation Firebird.
2000s–present
For the 2004 model year, the Pontiac GTO was relaunched in the United States as a rebadged captive import version of the Holden Monaro. The model was to recreate the past versions, but the new version "was nothing like the old aggressive and evocative model from the 60s" and it was discontinued in 2006.
For 2005, Chrysler introduced muscle car heritage to high-performance V8-powered versions of four-door sedans, the Dodge Charger and Chrysler 300C, using nameplates traditionally used for two-door muscle cars.
For 2005, the fifth-generation Ford Mustang, designed to resemble the original first-generation Mustang, brought back the aggressive lines and colors of the original.
For the 2006 model year, GM relaunched the Chevrolet Monte Carlo SS with the first V8 engine on the Monte Carlo in 15 Years. The same V8 was used on the Monte Carlo's W-Body sister cars like the Pontiac Grand Prix GXP, Buick Lacrosse Super, and the Chevrolet Impala SS. All Monte Carlo production ended on June 19, 2007, because of declining sales of coupe models in general as well as Chevrolet's plan to reintroduce a new Camaro.
For 2008, Chrysler re-introduced the Dodge Challenger, which features styling links to the 1970 first-generation Challenger and was claimed by the Chrysler CEO to be "a modern take on one of the most iconic muscle cars".
A year later, running on that same sentiment, Chevrolet released the 2009 fifth-generation Camaro, which bears some resemblance to the 1969 first-generation Camaro.
Australia
Origins
The first Australian-designed car to be marketed as a performance model was the 1963 Holden EH S4 model, of which 120 road cars were produced so that the model could be eligible to compete at the 1963 Armstrong 500 motor race at Bathurst. The EH S4 was powered by an upgraded version of the standard six-cylinder engine, enlarged to and producing . In 1964, the Ford Falcon (XM) became available with an enlarged "Super Pursuit" version of the standard six-cylinder engine, which produced .
In 1965, the Chrysler Valiant AP6 became the first Australian car to be available with a V8 engine. This optional engine was the version of the Chrysler LA engine, which produced and was imported from the United States. The first Australian-designed Ford to be available with a V8 was the 1966 Ford Falcon (XR), with a version of the Ford Windsor engine (imported from the United States), which produced . The first Holden to be available with a V8 was the 1968 Holden HK, with a version of the Chevrolet small-block V8 (imported from the United States) which produced . Later that year, a version of the engine became available in the Holden HK Monaro GTS 327 coupe.
The pinnacle of 1970s Australian muscle cars were the 1971–1972 Ford Falcon GTHO, Holden Monaro 350, and Chrysler Valiant Charger R/T (the smaller Holden Torana GTR was also a successful performance car of the era, but it is not considered a muscle car due to its prioritization of lighter weight over outright power output). The Ford Falcon (XY) GTHO Phase III model was powered by a version of the Ford Cleveland V8 engine, officially rated at , but estimated to produce between . The Holden HQ Monaro GTS 350 was powered by a version of the Chevrolet small-block V8 producing . The Chrysler Valiant Charger R/T E49 model was powered by a version of the Chrysler Hemi-6 six-cylinder engine producing .
Supercar scare
In 1972, the production of Australian muscle cars saw a setback when the Supercar scare caused Ford, Holden, and Chrysler to cease development of upcoming performance models, due to government pressure. The Australian muscle car models produced during the 1970s later consisted of the limited production 1977–1978 Holden Torana (LX) A9X option and the 1978–1979 Ford Falcon (XC) Cobra model, both created as homologation models for Group C touring car racing. These were less powerful than their predecessors.
Brands still offered high-performance models with V8 variants throughout the 1980s, but these vehicles were low production and were generally underpowered compared to their late 1960s and 1970s predecessors. An example was the Ford Falcon (XD), which was available with a 5.8 L V8 engine. Subsequent generations of the Ford Falcon would not have any V8 options available until 1992, when the EB XR8 was introduced. The Holden Commodore debuted in 1978. However, a renaissance in muscle cars would be sparked by factory-backed aftermarket operations. Holden Dealer Team would release high-performance models of the Holden Commodore throughout the 1980s, such as the HDT Group A, which would become iconic for its blue paintwork. In 1988, Ford released the Ford Falcon (EB), which was available with a V8 in a 25th anniversary special model celebrating the original Ford Falcon GT.
Argentina
In Argentina, Chevrolet and Dodge produced two acclaimed models of muscle cars. The first was the producer of the third American generation of the Chevrolet Nova, which in this country was renamed Chevy. This model was initially presented in a 4-door sedan version that maintained many physical features of the Nova coupe version, which would also be produced and marketed in Argentina a few years later. While Dodge produced in Argentina a series of vehicles based on the fourth generation of the Dodge Dart that received the name of Línea Dodge (Dodge Line). This vehicle presented sedan and coupe versions, which in turn were a local redesign of the Dart model and which, depending on its level of equipment, received different names (Polara, Coronado, RT, and GTX). In return for these brands, both Ford and the national producer IKA would respond with the production of two high-performance sedans, such as the Argentine version of the Ford Falcon and a derivative of the Rambler American model, called IKA Torino, which, in addition to its sedan version, would present a coupe version which would end up being acclaimed and popularized in the Argentine automotive field.
Lists of muscle cars (1962–1974)
According to Car and Driver, January 1990:
1964–1969 Pontiac GTO
1966–1971 Plymouth/Dodge A-body 426 models
1966–1967 Chevrolet Chevy II / Nova SS 327
1966–1969 Chevrolet Chevelle SS 396
1968–1969 Chevrolet Chevy II / Nova SS 396
1969 Ford Torino Cobra 428
1969 Plymouth Road Runner 440 Six Pack
1969 Dodge Super Bee 440 Six Pack
1969 Chevrolet Camaro ZL1
1970 Chevrolet Chevelle SS 454
According to CNBC, April 2013:
1968 Shelby Mustang GT500KR
1969 Ford Mustang Boss 429
1969 Dodge Charger Daytona Hemi
1969 Chevrolet Camaro ZL1
1970 Oldsmobile 442 W-30
1970 Buick GSX Stage 1
1970 Chevrolet Chevelle SS 454 LS6
1970 Pontiac GTO Judge Ram Air IV
1971 Plymouth Hemi 'Cuda
1974 Pontiac Firebird Trans Am SD455
According to Road & Track, January 2021:
1962 Pontiac Catalina Super Duty
1963 Studebaker Super Lark
1963 Chevrolet Impala Z11
1964 Ford Fairlane Thunderbolt
1967 Dodge Coronet W023
1968 Hurst Hemi Dart L023
1969 Chevrolet Camaro ZL1
1969–1970 Ford Mustang Boss 429
1970 Buick GSX Stage 1
1970 AMC "The Machine"
1970 Plymouth Hemi 'Cuda Super Track Pack
1971 Ford Mustang Mach 1 Drag Pack
According to Motor Trend, June 2023:
1962 Pontiac Catalina Super Duty
1963 Plymouth Savoy Max Wedge
1964 Ford Fairlane Thunderbolt
1965 Pontiac GTO Tri-Power
1966 Dodge Coronet Street Hemi
1968 AMC AMX
1969 Chevrolet Camaro COPO 427
1969 Ford Mustang Boss 429
1969 Dodge Charger Daytona Hemi
1970 Chevrolet Chevelle SS 454
| Technology | Motorized road transport | null |
228498 | https://en.wikipedia.org/wiki/Common%20brushtail%20possum | Common brushtail possum | The common brushtail possum (Trichosurus vulpecula, from the Greek for "furry tailed" and the Latin for "little fox", previously in the genus Phalangista) is a nocturnal, semiarboreal marsupial of the family Phalangeridae, native to Australia and invasive in New Zealand, and the second-largest of the possums.
Like most possums, the common brushtail possum is nocturnal. It is mainly a folivore, but has been known to eat small mammals such as rats. In most Australian habitats, eucalyptus leaves are a significant part of the diet, but rarely the sole item eaten. Its tail is prehensile and naked on its lower underside. The four colour variations are silver-grey, brown, black, and gold.
It is the Australian marsupial most often seen by city dwellers, as it is one of few that thrive in cities and a wide range of natural and human-modified environments. Around human habitations, common brushtails are inventive and determined foragers with a liking for fruit trees, vegetable gardens, and kitchen raids. Its once vast distribution has been greatly affected by drought, epizootic disease and intrusion of invasive mammals into its habitat.
The common brushtail possum was introduced to New Zealand in the 1850s to establish a fur industry, but in the mild subtropical climate of New Zealand, and with few to no natural predators, it thrived to the extent that it became a major agricultural and conservation pest.
Description
The common brushtail possum has large and pointed ears. Its bushy tail (hence its name) is adapted to grasping branches, prehensile at the end with a hairless ventral patch. Its fore feet have sharp claws and the first toe of each hind foot is clawless, but has a strong grasp. The possum grooms itself with the third and fourth toes which are fused together. It has a thick and woolly pelage that varies in colour depending on the subspecies. Colour patterns tend to be silver-grey, brown, black, red, or cream. The ventral areas are typically lighter and the tail is usually brown or black. The muzzle is marked with dark patches.
The common brushtail possum has a head and body length of with a tail length of . It weighs . Males are generally larger than females. In addition, the coat of the male tends to be reddish at the shoulders. As with most marsupials, the female brushtail possum has a forward-opening, well-developed pouch. The chest of both sexes has a scent gland that emits a reddish secretion which stains that fur around it. It marks its territory with these secretions.
Biology and ecology
Range and habitat
The common brushtail possum is perhaps the most widespread marsupial of Australia. It is found throughout the eastern and northern parts of the continent, as well as some western regions, Tasmania and a number of offshore islands, such as Kangaroo Island and Barrow Island. Western Australia alone has several scattered population groups locally distinguished with given indigenous names: nunguin in Kimberley, walambari in Pilbara, wayurta in the desert areas, and bilda in Nullarbor Plain shared with South Australia among many others.
It is also widespread in New Zealand since its introduction in 1850. The common brushtail possum can be found in a variety of habitats, such as forests, semi-arid areas and even cultivated or urban areas. It is mostly a forest inhabiting species, however it is also found in treeless areas. In New Zealand, possums favour broadleaf-podocarp near farmland pastures. In southern beech forests and pine plantations, possums are less common. Overall, brushtail possums are more densely populated in New Zealand than in their native Australia. This may be because Australia has more fragmented eucalypt forests and more predators. In Australia, brushtail possums are threatened by humans, tiger quolls, dogs, foxes, cats, goannas, carpet snakes, and powerful owls. In New Zealand, brushtail possums are threatened only by humans and cats. The IUCN highlight the population trend in Australia as decreasing.
The northern subspecies of the common brushtail possum has declined substantially, with one study in Australia's Northern Territory finding a 22% reduction in the extent of occurrence of and a 50% reduction in the breadth of occupied environmental space. Analysis of contemporary occurrence points suggested that the species is contracting towards areas of higher rainfall, lower fire frequency, and higher vegetation cover.
Food and foraging
The common brushtail possum can adapt to numerous kinds of vegetation but it is largely omnivorous. It prefers Eucalyptus leaves, but also eats flowers, shoots, fruits, and seeds. It may also consume animal matter such as insects, birds' eggs and chicks, and other small vertebrates. Brushtail possums may eat three or four different plant species during a foraging trip, unlike some other arboreal marsupials, such as the koala and the greater glider, which focus on single species. The brushtail possum's rounded molars cannot cut Eucalyptus leaves as finely as more specialised feeders. They are more adapted to crushing their food, which enables them to chew fruit or herbs more effectively. The brushtail possums' caecum lacks internal ridges and cannot separate coarse and fine particles as efficiently as some other arboreal marsupials. The brushtail possum cannot rely on Eucalyptus alone to provide sufficient protein. Its more generalised and mixed diet, however, does provide adequate nitrogen.
Behaviour
The common brushtail possum is largely arboreal and nocturnal. It has a mostly solitary lifestyle, and individuals keep their distance with scent markings (urinating) and vocalisations. They usually make their dens in natural places such as tree hollows and caves, but also use spaces in the roofs of houses. While they sometimes share dens, brushtails normally sleep in separate dens. Individuals from New Zealand use many more den sites than those from Australia. Brushtail possums compete with each other and other animals for den spaces, and this contributes to their mortality. This is likely another reason why brushtail possum population densities are smaller in Australia than in New Zealand. Brushtail possums are usually not aggressive towards each other and usually just stare with erect ears. They vocalise with clicks, grunts, hisses, alarm chatters, guttural coughs, and screeching.
Reproduction and life history
The common brushtail possum can breed at any time of the year, but breeding tends to peak in spring, from September to November, and in autumn, from March to May, in some areas. Mating is promiscuous and random; some males can sire several young in a season, while over half sire none. In one Queensland population, males apparently need a month of consorting with females before they can mate with them. Females have a gestation period of 16–18 days, after which they give birth to single young. A newborn brushtail possum is only long and weighs only . As usual for marsupials, the newborn may climb, unaided, through the female's fur and into the pouch and attach to a teat. The young develops and remains inside the mother's pouch for another 4–5 months. A preliminary study inducing ovulation through exposure of hormones resulted in changes to the appearance of mammary glands in females suggesting that mammary glands provide immunological protection to neonates through milk secretions. When older, the young is left in the den or rides on its mother's back until it is 7–9 months old. Females reach sexual maturity when they are a year old, and males do so at the end of their second year. Brushtail possums can live up to 13 years in the wild.
Female young have a higher survival rate than their male counterparts due to establishing their home ranges closer to their mothers, while males travel farther in search of new nesting sites, encountering established territories from which they may be forcibly ejected. In New Zealand's Ōrongorongo population, female young have been found to continue to associate with their mothers after weaning, and some inherit the prime den sites. A possible competition exists between mothers and daughters for dens, and daughters may be excluded from a den occupied by the mother. In forests with shortages of den sites, females apparently produce more sons, which do not compete directly for den sites, while in forests with plentiful den sites, female young are greater in number.
Relationship with humans
The common brushtail possum is considered a pest in some areas, as it is known to cause damage to pine plantations, regenerative forest, flowers, fruit trees, and buildings. Like other possums, it is rather tolerant of humans and can sometimes be hand fed, although it is not encouraged, as their claws are quite sharp and can cause infection or disease to humans if scratched. It is a traditional food source for some Indigenous Australian groups.
Australia
Although once hunted extensively for its fur, the common brushtail possum is largely protected throughout Australia. Tasmania gives crop-protection permits to landowners whose property has been damaged.
While its populations are declining in some regions due to habitat loss, urban populations indicate an adaptation to the presence of humans. In some mainland states, possum trapping is permitted when attempting to evict possums from human residences (e.g. roofs), but possums must be released after dusk within 24 hours of capture, no more than 50 m from the trapping site. In some states, e.g. Victoria, trapped possums may be taken to registered veterinarians to be killed. In South Australia, they are fully protected and permits are required for trapping possums in human residences or for keeping or rescuing sick or injured wild possums and other native animals. In Queensland, they can only be trapped by licensed commercial relocators who must release possums within of the point of capture to ensure that an animal is not released into another possum's home range (possums are less likely to survive if they are released into a new area where they do not have access to a den or must compete with a neighbouring possum for den space).
New Zealand
Since its introduction from Australia by European settlers in the 1850s, the common brushtail possum has become a major threat to New Zealand native forests and birds. It is also a host for the highly contagious bovine tuberculosis. This is not an issue in Australia, where the disease has been eradicated.
By the 1980s, the peak population had reached an estimated 60–70 million, but is now down to an estimated 30 million due to control measures. The New Zealand Department of Conservation controls possum numbers in many areas via the aerial dropping of 1080-laced bait. Hunting is not restricted, but the population seems to be stable despite the annual killing of thousands of the animals.
Hunting
Possums are a pest in New Zealand and Tasmania, where they are culled for their meat and fur. However, due to tuberculosis being prevalent in many possums across most of New Zealand, possums are generally only eaten in Northland, where the disease does not exist in possums. In Northland, possum meat has even been used in meat pies.
In Tasmania, possum meat is served at some restaurants. On Bruny Island, possum meat is sold at Bruny Island Game Meats, which also sell it at farmer's markets, including in Hobart.
| Biology and health sciences | Diprotodontia | Animals |
228613 | https://en.wikipedia.org/wiki/Starfish | Starfish | Starfish or sea stars are star-shaped echinoderms belonging to the class Asteroidea (). Common usage frequently finds these names being also applied to ophiuroids, which are correctly referred to as brittle stars or basket stars. Starfish are also known as asteroids due to being in the class Asteroidea. About 1,900 species of starfish live on the seabed in all the world's oceans, from warm, tropical zones to frigid, polar regions. They are found from the intertidal zone down to abyssal depths, at below the surface.
Starfish are marine invertebrates. They typically have a central disc and usually five arms, though some species have a larger number of arms. The aboral or upper surface may be smooth, granular or spiny, and is covered with overlapping plates. Many species are brightly coloured in various shades of red or orange, while others are blue, grey or brown. Starfish have tube feet operated by a hydraulic system and a mouth at the centre of the oral or lower surface. They are opportunistic feeders and are mostly predators on benthic invertebrates. Several species have specialized feeding behaviours including eversion of their stomachs and suspension feeding. They have complex life cycles and can reproduce both sexually and asexually. Most can regenerate damaged parts or lost arms and they can shed arms as a means of defense. The Asteroidea occupy several significant ecological roles. Starfish, such as the ochre sea star (Pisaster ochraceus) and the reef sea star (Stichaster australis), have become widely known as examples of the keystone species concept in ecology. The tropical crown-of-thorns starfish (Acanthaster planci) is a voracious predator of coral throughout the Indo-Pacific region, and the Northern Pacific seastar is on the list of the World's 100 Worst Invasive Alien Species.
The fossil record for starfish is ancient, dating back to the Ordovician around 450 million years ago, but it is rather sparse, as starfish tend to disintegrate after death. Only the ossicles and spines of the animal are likely to be preserved, making remains hard to locate. With their appealing symmetrical shape, starfish have played a part in literature, legend, design and popular culture. They are sometimes collected as curios, used in design or as logos, and in some cultures, despite possible toxicity, they are eaten.
Anatomy
Most starfish have five arms that radiate from a central disc, but the number varies with the group. Some species have six or seven arms and others have 10–15 arms. The Antarctic Labidiaster annulatus can have over fifty.
Mapping the expression patterns of genes that express differently across the body axes suggests that one could think of the body of a starfish as a disembodied head walking about the sea floor on its lips. The known markers for trunk structures are expressed only in internal tissues rather than on the surface. Only the front part of the axis, which specifies head-related structures, is represented on the body surface.
Body wall
The body wall consists of a thin cuticle, an epidermis consisting of a single layer of cells, a thick dermis formed of connective tissue and a thin coelomic myoepithelial layer, which provides the longitudinal and circular musculature. The dermis contains an endoskeleton of calcium carbonate components known as ossicles. These are honeycombed structures composed of calcite microcrystals arranged in a lattice. They vary in form, with some bearing external granules, tubercles and spines, but most are tabular plates that fit neatly together in a tessellated manner and form the main covering of the aboral surface. Some are specialised structures such as the madreporite (the entrance to the water vascular system), pedicellariae and paxillae. Pedicellariae are compound ossicles with forceps-like jaws. They remove debris from the body surface and wave around on flexible stalks in response to physical or chemical stimuli while continually making biting movements. They often form clusters surrounding spines. Paxillae are umbrella-like structures found on starfish that live buried in sediment. The edges of adjacent paxillae meet to form a false cuticle with a water cavity beneath in which the madreporite and delicate gill structures are protected. All the ossicles, including those projecting externally, are covered by the epidermal layer.
Several groups of starfish, including Valvatida and Forcipulatida, possess pedicellariae. In Forcipulatida, such as Asterias and Pisaster, they occur in pompom-like tufts at the base of each spine, whereas in the Goniasteridae, such as Hippasteria phrygiana, the pedicellariae are scattered over the body surface. Some are thought to assist in defence, while others aid in feeding or in the removal of organisms attempting to settle on the starfish's surface. Some species like Labidiaster annulatus, Rathbunaster californicus and Novodinia antillensis use their large pedicellariae to capture small fish and crustaceans.
There may also be papulae, thin-walled protrusions of the body cavity that reach through the body wall and extend into the surrounding water. These serve a respiratory function. The structures are supported by collagen fibres set at right angles to each other and arranged in a three-dimensional web with the ossicles and papulae in the interstices. This arrangement enables both easy flexion of the arms by the starfish and the rapid onset of stiffness and rigidity required for actions performed under stress.
Water vascular system
The water vascular system of the starfish is a hydraulic system made up of a network of fluid-filled canals and is concerned with locomotion, adhesion, food manipulation and gas exchange. Water enters the system through the madreporite, a porous, often conspicuous, sieve-like ossicle on the aboral surface. It is linked through a stone canal, often lined with calcareous material, to a ring canal around the mouth opening. A set of radial canals leads off this; one radial canal runs along the ambulacral groove in each arm. There are short lateral canals branching off alternately to either side of the radial canal, each ending in an ampulla. These bulb-shaped organs are joined to tube feet (podia) on the exterior of the animal by short linking canals that pass through ossicles in the ambulacral groove. There are usually two rows of tube feet but in some species, the lateral canals are alternately long and short and there appear to be four rows. The interior of the whole canal system is lined with cilia.
When longitudinal muscles in the ampullae contract, valves in the lateral canals close and water is forced into the tube feet. These extend to contact the substrate. Although the tube feet resemble suction cups in appearance, the gripping action is a function of adhesive chemicals rather than suction. Other chemicals and relaxation of the ampullae allow for release from the substrate. The tube feet latch on to surfaces and move in a wave, with one arm section attaching to the surface as another releases. Some starfish turn up the tips of their arms while moving which gives maximum exposure of the sensory tube feet and the eyespot to external stimuli.
Having descended from bilateral organisms, starfish may move in a bilateral fashion, particularly when hunting or in danger. When crawling, certain arms act as the leading arms, while others trail behind. Most starfish cannot move quickly, a typical speed being that of the leather star (Dermasterias imbricata), which can manage just in a minute. Some burrowing species from the genera Astropecten and Luidia have points rather than suckers on their long tube feet and are capable of much more rapid motion, "gliding" across the ocean floor. The sand star (Luidia foliolata) can travel at a speed of per minute. When a starfish finds itself upside down, two adjacent arms are bent backwards to provide support, the opposite arm is used to stamp the ground while the two remaining arms are raised on either side; finally the stamping arm is released as the starfish turns itself over and recovers its normal stance.
Apart from their function in locomotion, the tube feet act as accessory gills. The water vascular system serves to transport oxygen from, and carbon dioxide to, the tube feet and also nutrients from the gut to the muscles involved in locomotion. Fluid movement is bidirectional and initiated by cilia. Gas exchange also takes place through other gills known as papulae, which are thin-walled bulges on the aboral surface of the disc and arms. Oxygen is transferred from these to the coelomic fluid, which acts as the transport medium for gasses. Oxygen dissolved in the water is distributed through the body mainly by the fluid in the main body cavity; the circulatory system may also play a minor role.
Digestive system and excretion
The gut of a starfish occupies most of the disc and extends into the arms. The mouth is located in the centre of the oral surface, where it is surrounded by a tough peristomial membrane and closed with a sphincter. The mouth opens through a short oesophagus into a stomach divided by a constriction into a larger, eversible cardiac portion and a smaller pyloric portion. The cardiac stomach is glandular and pouched, and is supported by ligaments attached to ossicles in the arms so it can be pulled back into position after it has been everted. The pyloric stomach has two extensions into each arm: the pyloric caeca. These are elongated, branched hollow tubes that are lined by a series of glands, which secrete digestive enzymes and absorb nutrients from the food. A short intestine and rectum run from the pyloric stomach to open at a small anus at the apex of the aboral surface of the disc.
Primitive starfish, such as Astropecten and Luidia, swallow their prey whole, and start to digest it in their cardiac stomachs. Shell valves and other inedible materials are ejected through their mouths. The semi-digested fluid is passed into their pyloric stomachs and caeca where digestion continues and absorption ensues. In more advanced species of starfish, the cardiac stomach can be everted from the organism's body to engulf and digest food. When the prey is a clam or other bivalve, the starfish pulls with its tube feet to separate the two valves slightly, and inserts a small section of its stomach, which releases enzymes to digest the prey. The stomach and the partially digested prey are later retracted into the disc. Here the food is passed on to the pyloric stomach, which always remains inside the disc. The retraction and contraction of the cardiac stomach is activated by a neuropeptide known as NGFFYamide.
Because of this ability to digest food outside the body, starfish can hunt prey much larger than their mouths. Their diets include clams and oysters, arthropods, small fish and gastropod molluscs. Some starfish are not pure carnivores, supplementing their diets with algae or organic detritus. Some of these species are grazers, but others trap food particles from the water in sticky mucus strands that are swept towards the mouth along ciliated grooves.
The main nitrogenous waste product is ammonia. Starfish have no distinct excretory organs; waste ammonia is removed by diffusion through the tube feet and papulae. The body fluid contains phagocytic cells called coelomocytes, which are also found within the hemal and water vascular systems. These cells engulf waste material, and eventually migrate to the tips of the papulae, where a portion of body wall is nipped off and ejected into the surrounding water. Some waste may also be excreted by the pyloric glands and voided with the faeces.
Starfish do not appear to have any mechanisms for osmoregulation, and keep their body fluids at the same salt concentration as the surrounding water. Although some species can tolerate relatively low salinity, the lack of an osmoregulation system probably explains why starfish are not found in fresh water or even in many estuarine environments.
Sensory and nervous systems
Although starfish do not have many well-defined sense organs, they are sensitive to touch, light, temperature, orientation and the status of the water around them. The tube feet, spines and pedicellariae are sensitive to touch. The tube feet, especially those at the tips of the rays, are also sensitive to chemicals, enabling the starfish to detect odour sources such as food. There are eyespots at the ends of the arms, each one made of 80–200 simple ocelli. These are composed of pigmented epithelial cells that respond to light and are covered by a thick, transparent cuticle that both protects the ocelli and acts to focus light. Many starfish also possess individual photoreceptor cells in other parts of their bodies and respond to light even when their eyespots are covered. Whether they advance or retreat depends on the species.
While a starfish lacks a centralized brain, it has a complex nervous system with a nerve ring around the mouth and a radial nerve running along the ambulacral region of each arm parallel to the radial canal. The peripheral nerve system consists of two nerve nets: a sensory system in the epidermis and a motor system in the lining of the coelomic cavity. Neurons passing through the dermis connect the two. The ring nerves and radial nerves have sensory and motor components and coordinate the starfish's balance and directional systems. The sensory component receives input from the sensory organs while the motor nerves control the tube feet and musculature. The starfish does not have the capacity to plan its actions. If one arm detects an attractive odour, it becomes dominant and temporarily over-rides the other arms to initiate movement towards the prey. The mechanism for this is not fully understood.
Circulatory system
The body cavity contains the circulatory or haemal system. The vessels form three rings: one around the mouth (the hyponeural haemal ring), another around the digestive system (the gastric ring) and the third near the aboral surface (the genital ring). The heart beats about six times a minute and is at the apex of a vertical channel (the axial vessel) that connects the three rings. At the base of each arm are paired gonads; a lateral vessel extends from the genital ring past the gonads to the tip of the arm. This vessel has a blind end and there is no continuous circulation of the fluid within it. This liquid does not contain a pigment and has little or no respiratory function but is probably used to transport nutrients around the body.
Secondary metabolites
Starfish produce a large number of secondary metabolites in the form of lipids, including steroidal derivatives of cholesterol, and fatty acid amides of sphingosine. The steroids are mostly saponins, known as asterosaponins, and their sulphated derivatives. They vary between species and are typically formed from up to six sugar molecules (usually glucose and galactose) connected by up to three glycosidic chains. Long-chain fatty acid amides of sphingosine occur frequently and some of them have known pharmacological activity. Various ceramides are also known from starfish and a small number of alkaloids have also been identified. The functions of these chemicals in the starfish have not been fully investigated but most have roles in defence and communication. Some are feeding deterrents used by the starfish to discourage predation. Others are antifoulants and supplement the pedicellariae to prevent other organisms from settling on the starfish's aboral surface. Some are alarm pheromones and escape-eliciting chemicals, the release of which trigger responses in conspecific starfish but often produce escape responses in potential prey. Research into the efficacy of these compounds for possible pharmacological or industrial use occurs worldwide.
Life cycle
Sexual reproduction
Most species of starfish are gonochorous, there being separate male and female individuals. These are usually not distinguishable externally as the gonads cannot be seen, but their sex is apparent when they spawn. Some species are simultaneous hermaphrodites, producing eggs and sperm at the same time, and in a few of these the same gonad, called an ovotestis, produces both eggs and sperm. Other starfish are sequential hermaphrodites. Protandrous individuals of species like Asterina gibbosa start life as males before changing sex into females as they grow older. In some species such as Nepanthia belcheri, a large female can split in half and the resulting offspring are males. When these grow large enough they change back into females.
Each starfish arm contains two gonads that release gametes through openings called gonoducts, located on the central disc between the arms. Fertilization is generally external but in a few species, internal fertilization takes place. In most species, the buoyant eggs and sperm are simply released into the water (free spawning) and the resulting embryos and larvae live as part of the plankton. In others, the eggs may be stuck to the undersides of rocks. In certain species of starfish, the females brood their eggs – either by simply enveloping them or by holding them in specialised structures. Brooding may be done in pockets on the starfish's aboral surface, inside the pyloric stomach (Leptasterias tenera) or even in the interior of the gonads themselves. Those starfish that brood their eggs by "sitting" on them usually assume a humped posture with their discs raised off the substrate. Pteraster militaris broods a few of its young and disperses the remaining eggs, that are too numerous to fit into its pouch. In these brooding species, the eggs are relatively large, and supplied with yolk, and they generally develop directly into miniature starfish without an intervening larval stage. The developing young are called lecithotrophic because they obtain their nutrition from the yolk as opposed to "planktotrophic" larvae that feed in the water column. In Parvulastra parvivipara, an intragonadal brooder, the young starfish obtain nutrients by eating other eggs and embryos in the brood pouch. Brooding is especially common in polar and deep-sea species that live in environments unfavourable for larval development and in smaller species that produce just a few eggs.
In the tropics, a plentiful supply of phytoplankton is continuously available for starfish larvae to feed on. Spawning takes place at any time of year, each species having its own characteristic breeding season. In temperate regions, the spring and summer brings an increase in food supplies. The first individual of a species to spawn may release a pheromone that serves to attract other starfish to aggregate and to release their gametes synchronously. In other species, a male and female may come together and form a pair. This behaviour is called pseudocopulation and the male climbs on top, placing his arms between those of the female. When she releases eggs into the water, he is induced to spawn. Starfish may use environmental signals to coordinate the time of spawning (day length to indicate the correct time of the year, dawn or dusk to indicate the correct time of day), and chemical signals to indicate their readiness to breed. In some species, mature females produce chemicals to attract sperm in the sea water.
Larval development
Most starfish embryos hatch at the blastula stage. The original ball of cells develops a lateral pouch, the archenteron. The entrance to this is known as the blastopore and it will later develop into the anus—together with chordates, echinoderms are deuterostomes, meaning the second (deutero) invagination becomes the mouth (stome); members of all other phyla are protostomes, and their first invagination becomes the mouth. Another invagination of the surface will fuse with the tip of the archenteron as the mouth while the interior section will become the gut. At the same time, a band of cilia develops on the exterior. This enlarges and extends around the surface and eventually onto two developing arm-like outgrowths. At this stage the larva is known as a bipinnaria. The cilia are used for locomotion and feeding, their rhythmic beat wafting phytoplankton towards the mouth.
The next stage in development is a brachiolaria larva and involves the growth of three short, additional arms. These are at the anterior end, surround a sucker and have adhesive cells at their tips. Both bipinnaria and brachiolaria larvae are bilaterally symmetrical. When fully developed, the brachiolaria settles on the seabed and attaches itself with a short stalk formed from the ventral arms and sucker. Metamorphosis now takes place with a radical rearrangement of tissues. The left side of the larval body becomes the oral surface of the juvenile and the right side the aboral surface. Part of the gut is retained, but the mouth and anus move to new positions. Some of the body cavities degenerate but others become the water vascular system and the visceral coelom. The starfish is now pentaradially symmetrical. It casts off its stalk and becomes a free-living juvenile starfish about in diameter. Starfish of the order Paxillosida have no brachiolaria stage, with the bipinnaria larvae settling on the seabed and developing directly into juveniles.
Asexual reproduction
Some species of starfish in the three families Asterinidae, Asteriidae and Solasteridae are able to reproduce asexually as adults either by fission of their central discs or by autotomy of one or more of their arms. Which of these processes occurs depends on the genus. Among starfish that are able to regenerate their whole body from a single arm, some can do so even from fragments just long. Single arms that regenerate a whole individual are called comet forms. The division of the starfish, either across its disc or at the base of the arm, is usually accompanied by a weakness in the structure that provides a fracture zone.
The larvae of several species of starfish can reproduce asexually before they reach maturity. They do this by autotomising some parts of their bodies or by budding. When such a larva senses that food is plentiful, it takes the path of asexual reproduction rather than normal development. Though this costs it time and energy and delays maturity, it allows a single larva to give rise to multiple adults when the conditions are appropriate.
Regeneration
Some species of starfish have the ability to regenerate lost arms and can regrow an entire new limb given time. A few can regrow a complete new disc from a single arm, while others need at least part of the central disc to be attached to the detached part. Regrowth can take several months or years, and starfish are vulnerable to infections during the early stages after the loss of an arm. A separated limb lives off stored nutrients until it regrows a disc and mouth and is able to feed again. Other than fragmentation carried out for the purpose of reproduction, the division of the body may happen inadvertently due to part being detached by a predator, or part may be actively shed by the starfish in an escape response. The loss of parts of the body is achieved by the rapid softening of a special type of connective tissue in response to nervous signals. This type of tissue is called catch connective tissue and is found in most echinoderms. An autotomy-promoting factor has been identified which, when injected into another starfish, causes rapid shedding of arms.
Lifespan
The lifespan of a starfish varies considerably between species, generally being longer in larger forms and in those with planktonic larvae. For example, Leptasterias hexactis broods a small number of large-yolked eggs. It has an adult weight of , reaches sexual maturity in two years and lives for about ten years. Pisaster ochraceus releases a large number of eggs into the sea each year and has an adult weight of up to . It reaches maturity in five years and has a maximum recorded lifespan of 34 years. The average lifespan of a starfish is 35 years, and larger starfish species typically live longer than their smaller counterparts.
Ecology
Distribution and habitat
Echinoderms, including starfish, maintain a delicate internal electrolyte balance that is in equilibrium with sea water, making it impossible for them to live in a freshwater habitat. Starfish species inhabit all of the world's oceans. Habitats range from tropical coral reefs, rocky shores, tidal pools, mud, and sand to kelp forests, seagrass meadows and the deep-sea floor down to at least . The greatest diversity of species occurs in coastal areas.
Diet
Most species are generalist predators, eating microalgae, sponges, bivalves, snails and other small animals. The crown-of-thorns starfish consumes coral polyps, while other species are detritivores, feeding on decomposing organic material and faecal matter. A few are suspension feeders, gathering in phytoplankton; Henricia and Echinaster often occur in association with sponges, benefiting from the water current they produce. Various species have been shown to be able to absorb organic nutrients from the surrounding water, and this may form a significant portion of their diet.
The processes of feeding and capture may be aided by special parts; Pisaster brevispinus, the short-spined pisaster from the West Coast of America, can use a set of specialized tube feet to dig itself deep into the soft substrate to extract prey (usually clams). Grasping the shellfish, the starfish slowly pries open the prey's shell by wearing out its adductor muscle, and then inserts its everted stomach into the crack to digest the soft tissues. The gap between the valves need only be a fraction of a millimetre wide for the stomach to gain entry. Cannibalism has been observed in juvenile sea stars as early as four days after metamorphosis.
Ecological impact
Starfish are keystone species in their respective marine communities. Their relatively large sizes, diverse diets and ability to adapt to different environments makes them ecologically important. The term "keystone species" was in fact first used by Robert Paine in 1966 to describe a starfish, Pisaster ochraceus. When studying the low intertidal coasts of Washington state, Paine found that predation by P. ochraceus was a major factor in the diversity of species. Experimental removals of this top predator from a stretch of shoreline resulted in lower species diversity and the eventual domination of Mytilus mussels, which were able to outcompete other organisms for space and resources. Similar results were found in a 1971 study of Stichaster australis on the intertidal coast of the South Island of New Zealand. S. australis was found to have removed most of a batch of transplanted mussels within two or three months of their placement, while in an area from which S. australis had been removed, the mussels increased in number dramatically, overwhelming the area and threatening biodiversity.
The feeding activity of the omnivorous starfish Oreaster reticulatus on sandy and seagrass bottoms in the Virgin Islands appears to regulate the diversity, distribution and abundance of microorganisms. These starfish engulf piles of sediment removing the surface films and algae adhering to the particles. Organisms that dislike this disturbance are replaced by others better able to rapidly recolonise "clean" sediment. In addition, foraging by these migratory starfish creates diverse patches of organic matter, which may play a role in the distribution and abundance of organisms such as fish, crabs and sea urchins that feed on the sediment.
Starfish sometimes have negative effects on ecosystems. Outbreaks of crown-of-thorns starfish have caused damage to coral reefs in Northeast Australia and French Polynesia. A study in Polynesia found that coral cover declined drastically with the arrival of migratory starfish in 2006, dropping from 50% to under 5% in three years. This had a cascading effect on the whole benthic community and reef-feeding fish. Asterias amurensis is one of a few echinoderm invasive species. Its larvae likely arrived in Tasmania from central Japan via water discharged from ships in the 1980s. The species has since grown in numbers to the point where they threaten commercially important bivalve populations. As such, they are considered pests, and are on the Invasive Species Specialist Group's list of the world's 100 worst invasive species.
Sea Stars (starfish) are the main predators of kelp-eating sea urchins. Satellite imagery shows that sea urchin populations have exploded due to starfish mass deaths, and that by 2021, sea urchins have destroyed 95% of California's kelp forests.
Threats
Starfish may be preyed on by conspecifics, sea anemones, other starfish species, tritons, crabs, fish, gulls and sea otters. Their first lines of defence are the saponins present in their body walls, which have unpleasant flavours. Some starfish such as Astropecten polyacanthus also include powerful toxins such as tetrodotoxin among their chemical armoury, and the slime star can ooze out large quantities of repellent mucus. They also have body armour in the form of hard plates and spines. The crown-of-thorns starfish is particularly unattractive to potential predators, being heavily defended by sharp spines, laced with toxins and sometimes with bright warning colours. Other species protect their vulnerable tube feet and arm tips by lining their ambulacral grooves with spines and heavily plating their extremities.
Several species sometimes suffer from a wasting condition caused by bacteria in the genus Vibrio; however, a more widespread wasting disease, causing mass mortalities among starfish, appears sporadically. A paper published in November 2014 revealed the most likely cause of this disease to be a densovirus the authors named sea star-associated densovirus (SSaDV).
The protozoan Orchitophrya stellarum is known to infect the gonads of starfish and damage tissue. Starfish are vulnerable to high temperatures. Experiments have shown that the feeding and growth rates of P. ochraceus reduce greatly when their body temperatures rise above and that they die when their temperature rises to . This species has a unique ability to absorb seawater to keep itself cool when it is exposed to sunlight by a receding tide. It also appears to rely on its arms to absorb heat, so as to protect the central disc and vital organs like the stomach.
Starfish and other echinoderms are sensitive to marine pollution. The common starfish is considered to be a bioindicator for marine ecosystems. A 2009 study found that P. ochraceus is unlikely to be affected by ocean acidification as severely as other marine animals with calcareous skeletons. In other groups, structures made of calcium carbonate are vulnerable to dissolution when the pH is lowered. Researchers found that when P. ochraceus were exposed to and 770 ppm carbon dioxide (beyond rises expected in the next century), they were relatively unaffected. Their survival is likely due to the nodular nature of their skeletons, which are able to compensate for a shortage of carbonate by growing more fleshy tissue.
Evolution
Fossil record
Echinoderms first appeared in the fossil record in the Cambrian. The first known asterozoans were the Somasteroidea, which exhibit characteristics of both groups. Starfish are infrequently found as fossils, possibly because their hard skeletal components separate as the animal decays. Despite this, there are a few places where accumulations of complete skeletal structures occur, fossilized in place in Lagerstätten – so-called "starfish beds".
By the late Paleozoic, the crinoids and blastoids were the predominant echinoderms, and some limestones from this period are made almost entirely from fragments from these groups. In the two major extinction events that occurred during the late Devonian and late Permian, the blastoids were wiped out and only a few species of crinoids survived. Many starfish species also became extinct in these events, but afterwards the surviving few species diversified rapidly within about sixty million years during the Early Jurassic and the beginning of the Middle Jurassic. A 2012 study found that speciation in starfish can occur rapidly. During the last 6,000 years, divergence in the larval development of Cryptasterina hystera and Cryptasterina pentagona has taken place, the former adopting internal fertilization and brooding and the latter remaining a broadcast spawner.
Diversity
The scientific name Asteroidea was given to starfish by the French zoologist de Blainville in 1830. It is derived from the Greek aster, ἀστήρ (a star) and the Greek eidos, εἶδος (form, likeness, appearance). The class Asteroidea belongs to the phylum Echinodermata. As well as the starfish, the echinoderms include sea urchins, sand dollars, brittle and basket stars, sea cucumbers and crinoids. The larvae of echinoderms have bilateral symmetry, but during metamorphosis this is replaced with radial symmetry, typically pentameric. Adult echinoderms are characterized by having a water vascular system with external tube feet and a calcareous endoskeleton consisting of ossicles connected by a mesh of collagen fibres. Starfish are included in the subphylum Asterozoa, the characteristics of which include a flattened, star-shaped body as adults consisting of a central disc and multiple radiating arms. The subphylum includes the two classes of Asteroidea, the starfish, and Ophiuroidea, the brittle stars and basket stars. Asteroids have broad-based arms with skeletal support provided by calcareous plates in the body wall while ophiuroids have clearly demarcated slender arms strengthened by paired fused ossicles forming jointed "vertebrae".
The starfish are a large and diverse class with over 1,900 living species. There are seven extant orders, Brisingida, Forcipulatida, Notomyotida, Paxillosida, Spinulosida, Valvatida and Velatida and two extinct ones, Calliasterellidae and Trichasteropsida. Living asteroids, the Neoasteroidea, are morphologically distinct from their forerunners in the Paleozoic. The taxonomy of the group is relatively stable but there is ongoing debate about the status of the Paxillosida, and the deep-water sea daisies, though clearly Asteroidea and currently included in Velatida, do not fit easily in any accepted lineage. Phylogenetic data suggests that they may be a sister group, the Concentricycloidea, to the Neoasteroidea, or that the Velatida themselves may be a sister group.
Living groups
Brisingida (2 families, 17 genera, 111 species)
Species in this order have a small, inflexible disc and 6–20 long, thin arms, which they use for suspension feeding. They have a single series of marginal plates, a fused ring of disc plates, a reduced number of aboral plates, crossed pedicellariae, and several series of long spines on the arms. They live almost exclusively in deep-sea habitats, although a few live in shallow waters in the Antarctic. In some species, the tube feet have rounded tips and lack suckers.
Forcipulatida (6 families, 63 genera, 269 species)
Species in this order have distinctive pedicellariae, consisting of a short stalk with three skeletal ossicles. They tend to have robust bodies and have tube feet with flat-tipped suckers usually arranged in four rows. The order includes well-known species from temperate regions, including the common starfish of North Atlantic coasts and rock pools, as well as cold-water and abyssal species.
Notomyotida (1 family, 8 genera, 75 species)
These starfish are deep-sea dwelling and have particularly flexible arms. The inner dorso-lateral surfaces of the arms contain characteristic longitudinal muscle bands. In some species, the tube feet lack suckers.
Paxillosida (7 families, 48 genera, 372 species)
This is a primitive order and members do not extrude their stomach when feeding, lack an anus and have no suckers on their tube feet. Papulae are plentiful on their aboral surface and they possess marginal plates and paxillae. They mostly inhabit soft-bottomed areas of sand or mud. There is no brachiolaria stage in their larval development. The comb starfish (Astropecten polyacanthus) is a member of this order.
Spinulosida (1 family, 8 genera, 121 species)
Most species in this order lack pedicellariae and all have a delicate skeletal arrangement with small or no marginal plates on the disc and arms. They have numerous groups of short spines on the aboral surface. This group includes the red starfish Echinaster sepositus.
Valvatida (16 families, 172 genera, 695 species)
Most species in this order have five arms and two rows of tube feet with suckers. There are conspicuous marginal plates on the arms and disc. Some species have paxillae and in some, the main pedicellariae are clamp-like and recessed into the skeletal plates. This group includes the cushion stars, the leather star and the sea daisies.
Velatida (4 families, 16 genera, 138 species)
This order of starfish consists mostly of deep-sea and other cold-water starfish often with a global distribution. The shape is pentagonal or star-shaped with five to fifteen arms. They mostly have poorly developed skeletons with papulae widely distributed on the aboral surface and often spiny pedicellariae. This group includes the slime star.
Extinct groups
Extinct groups within the Asteroidea include:
† Calliasterellidae, with the type genus Calliasterella from the Devonian and Carboniferous
† Palastericus, a Devonian genus
† Trichasteropsida, with the Triassic genus Trichasteropsis (at least 2 species)
Phylogeny
External
Starfish are deuterostome animals, like the chordates. A 2014 analysis of 219 genes from all classes of echinoderms gives the following phylogenetic tree. The times at which the clades diverged are shown under the labels in millions of years ago (mya).
Internal
The phylogeny of the Asteroidea has been difficult to resolve, with visible (morphological) features proving inadequate, and the question of whether traditional taxa are clades in doubt. The phylogeny proposed by Gale in 1987 is:
The phylogeny proposed by Blake in 1987 is:
Later work making use of molecular evidence, with or without the use of morphological evidence, had by 2000 failed to resolve the argument. In 2011, on further molecular evidence, Janies and colleagues noted that the phylogeny of the echinoderms "has proven difficult", and that "the overall phylogeny of extant echinoderms remains sensitive to the choice of analytical methods". They presented a phylogenetic tree for the living Asteroidea only; using the traditional names of starfish orders where possible, and indicating "part of" otherwise, the phylogeny is shown below. The Solasteridae are split from the Velatida, and the old Spinulosida is broken up.
Human relations
In research
Starfish are deuterostomes, closely related, together with all other echinoderms, to chordates, and are used in reproductive and developmental studies. Female starfish produce large numbers of oocytes that are easily isolated; these can be stored in a pre-meiosis phase and stimulated to complete division by the use of 1-methyladenine. Starfish oocytes are well suited for this research as they are large and easy to handle, transparent, simple to maintain in sea water at room temperature, and they develop rapidly. Asterina pectinifera, used as a model organism for this purpose, is resilient and easy to breed and maintain in the laboratory.
Another area of research is the ability of starfish to regenerate lost body parts. The stem cells of adult humans are incapable of much differentiation and understanding the regrowth, repair and cloning processes in starfish may have implications for human medicine.
Starfish also have an unusual ability to expel foreign objects from their bodies, which makes them difficult to tag for research tracking purposes.
In legend and culture
An aboriginal Australian fable retold by the Welsh school headmaster William Jenkyn Thomas (1870–1959) tells how some animals needed a canoe to cross the ocean. Whale had one but refused to lend it, so Starfish kept him busy, telling him stories and grooming him to remove parasites, while the others stole the canoe. When Whale realized the trick he beat Starfish ragged, which is how Starfish still is today.
In 1900, the scholar Edward Tregear documented The Creation Song, which he describes as "an ancient prayer for the dedication of a high chief" of Hawaii. Among the "uncreated gods" described early in the song are the male Kumulipo ("Creation") and the female Poele, both born in the night, a coral insect, the earthworm, and the starfish.
Georg Eberhard Rumpf's 1705 The Ambonese Curiosity Cabinet describes the tropical varieties of Stella Marina or Bintang Laut, "Sea Star", in Latin and Malay respectively, known in the waters around Ambon. He writes that the Histoire des Antilles reports that when the sea stars "see thunder storms approaching, [they] grab hold of many small stones with their little legs, looking to ... hold themselves down as if with anchors".
Starfish is the title of novels by Peter Watts and Jennie Orbell, and in 2012, Alice Addison wrote a non-fiction book titled Starfish: A Year in the Life of Bereavement and Depression. The Starfish and the Spider is a 2006 business management book by Ori Brafman and Rod Beckstrom; its title alludes to the ability of the starfish to regenerate itself because of its decentralized nervous system, and the book suggests ways that a decentralized organisation may flourish.
In the Nickelodeon animated television series SpongeBob SquarePants, the eponymous character's best friend is a dim-witted starfish, Patrick Star.
As food
Starfish are widespread in the oceans, but are only occasionally used as food. There may be good reason for this: the bodies of numerous species are dominated by bony ossicles, and the body wall of many species contains saponins, which have an unpleasant taste, and others contain tetrodotoxins which are poisonous. Some species that prey on bivalve molluscs can transmit paralytic shellfish poisoning. Georg Eberhard Rumpf found few starfish being used for food in the Indonesian archipelago, other than as bait in fish traps, but on the island of "Huamobel" the people cut them up, squeeze out the "black blood" and cook them with sour tamarind leaves; after resting the pieces for a day or two, they remove the outer skin and cook them in coconut milk. Starfish are sometimes eaten in China, Japan and in Micronesia.
As collectables
Starfish are in some cases taken from their habitat and sold to tourists as souvenirs, ornaments, curios or for display in aquariums. In particular, Oreaster reticulatus, with its easily accessed habitat and conspicuous coloration, is widely collected in the Caribbean. In the early to mid 20th century, this species was common along the coasts of the West Indies, but collection and trade have severely reduced its numbers. In the State of Florida, O. reticulatus is listed as endangered and its collection is illegal. Nevertheless, it is still sold throughout its range and beyond. A similar phenomenon exists in the Indo-Pacific for species such as Protoreaster nodosus.
In industry and military history
With its multiple arms, the starfish provides a popular metaphor for computer networks, companies and software tools. It is also the name of a seabed imaging system and company.
Starfish has repeatedly been chosen as a name in military history. Three ships of the Royal Navy have borne the name HMS Starfish: an A-class destroyer launched in 1894; an R-class destroyer launched in 1916; and an S-class submarine launched in 1933 and lost in 1940. In World War II, Starfish sites were large-scale night-time decoys created during The Blitz to simulate burning British cities. Starfish Prime was a high-altitude nuclear test conducted by the United States on 9 July 1962.
| Biology and health sciences | Echinoderms | null |
228744 | https://en.wikipedia.org/wiki/Bottle | Bottle | A bottle is a narrow-necked container made of an impermeable material (such as glass, plastic or aluminium) in various shapes and sizes that stores and transports liquids. Its mouth, at the bottling line, can be sealed with an internal stopper, an external bottle cap, a closure, or induction sealing.
Etymology
First attested in 14th century. From the English word bottle derives from an Old French word boteille, from vulgar Latin butticula, from late Latin buttis ("cask"), a latinisation of the Greek βοῦττις (bouttis) ("vessel").
Types
Glass
Wine
The glass bottle represented an important development in the history of wine, because, when combined with a high-quality stopper such as a cork, it allowed long-term aging of wine. Glass has all the qualities required for long-term storage. It eventually gave rise to "château bottling", the practice where an estate's wine is put in a bottle at the source, rather than by a merchant. Prior to this, wine used to be sold by the barrel (and before that, the amphora) and put into bottles only at the merchant's shop, if at all. This left large and often abused opportunities for fraud and adulteration, as consumers had to trust the merchant as to the contents. It is thought that most wine consumed outside of wine-producing regions had been tampered with in some way. Also, not all merchants were careful to avoid oxidation or contamination while bottling, leading to large bottle variation. Particularly in the case of port, certain conscientious merchants' bottling of old ports fetch higher prices even today. To avoid these problems, most fine wine is bottled at the place of production (including all port, since 1974).
There are many sizes and shapes of bottles used for wine. Some of the known shapes:
"Bordeaux": This bottle is roughly straight sided with a curved "shoulder" that is useful for catching sediment and is also the easiest to stack. Traditionally used in Bordeaux but now worldwide, this is probably the most common type.
"Burgundy": Traditionally used in Burgundy, this has sides that taper down about 2/3 of the height to a short cylindrical section, and does not have a shoulder.
"Champagne": Traditionally used for Champagne, it is similar to a Burgundy bottle, but with a wider base and heavier construction to withstand the pressure from the carbonation of the sparkling wine.
Codd-neck
In 1872, British soft drink makers Hiram Codd of Camberwell, London, designed and patented a bottle designed specifically for carbonated drinks. The Codd-neck bottle was designed and manufactured to enclose a marble and a rubber washer/gasket in the neck. The bottles were filled upside down, and pressure of the gas in the bottle forced the marble against the washer, sealing in the carbonation. The bottle was pinched into a special shape, as can be seen in the photo to the left, to provide a chamber into which the marble was pushed to open the bottle. This prevented the marble from blocking the neck as the drink was poured.
Soon after its introduction, the bottle became extremely popular with the soft drink and brewing industries, mainly in Europe, Asia and Australasia, though some alcohol drinkers disdained the use of the bottle. One etymology of the term codswallop originates from beer sold in Codd bottles, though this is generally dismissed as a folk etymology.
The bottles were regularly produced for many decades, but gradually declined in usage. Since children smashed the bottles to retrieve the marbles, they are relatively scarce and have become collector items; particularly in the UK. A cobalt-coloured Codd bottle today fetches hundreds of British pounds at auction. The Codd-neck design is still used for the Japanese soft drink Ramune and in the Indian drink called Banta.
Plastic
The plastic is strain oriented in the stretch blow molding manufacturing process. Plastic bottles are typically used to store liquids such as water, soft drinks, motor oil, cooking oil, medicine, shampoo, milk, and ink. The size ranges from very small sample bottles to very large carboys. The main advantages of plastic bottles over glass are their superior resistance to breakage, in both production and transportation, as well as their light weight and low cost of production. Disadvantages include widespread plastic pollution.
Aluminium
An aluminium bottle is a bottle made of aluminium (or aluminum, outside of British English). In some countries, it is also called a "bottlecan". It usually holds beer, soft drinks or wine.
Hot water
A hot water bottle is a bottle filled with hot water used to provide warmth. It can be made from various materials, most commonly rubber, but has historically been made from harder materials such as metal, glass, earthenware, or wood.
Gallery
Miscellany
Bottles are often recycled according to the SPI recycling code for the material.
| Technology | Containers | null |
228798 | https://en.wikipedia.org/wiki/Iguanodon | Iguanodon | Iguanodon ( ; meaning 'iguana-tooth'), named in 1825, is a genus of iguanodontian dinosaur. While many species found worldwide have been classified in the genus Iguanodon, dating from the Late Jurassic to Early Cretaceous, taxonomic revision in the early 21st century has defined Iguanodon to be based on one well-substantiated species: I. bernissartensis, which lived during the Barremian to early Aptian ages of the Early Cretaceous in Belgium, Germany, England, and Spain, between about 126 and 122 million years ago. Iguanodon was a large, bulky herbivore, measuring up to in length and in body mass. Distinctive features include large thumb spikes, which were possibly used for defense against predators, combined with long prehensile fifth fingers able to forage for food.
The genus was named in 1825 by English geologist Gideon Mantell, based on fossil specimens found in England and was given the species name I. anglicus. Iguanodon was the second type of dinosaur formally named based on fossil specimens, after Megalosaurus. Together with Megalosaurus and Hylaeosaurus, it was one of the three genera originally used to define Dinosauria. The genus Iguanodon belongs to the larger group Iguanodontia, along with the duck-billed hadrosaurs. The taxonomy of this genus continues to be a topic of study as new species are named or long-standing ones reassigned to other genera.
In 1878 new, far more complete remains of Iguanodon were discovered in Belgium and studied by Louis Dollo. These were given the new species I. bernissartensis. In the early 21st century it became understood that the remains referred to as Iguanodon in England belonged to four different species (including I. bernissartensis) that were not closely related to each other, which were subsequently split off into Mantellisaurus, Barilium and Hypselospinus. It was also found that the originally described type species of Iguanodon, I. anglicus is now a nomen dubium, and not valid. Thus the name "Iguanodon" became fixed around the well known species based primarily on the Belgian specimens. In 2015, a second valid species, I. galvensis, was named, based on fossils found in the Iberian Peninsula.
Scientific understanding of Iguanodon has evolved over time as new information has been obtained from fossils. The numerous specimens of this genus, including nearly complete skeletons from two well-known bone beds, have allowed researchers to make informed hypotheses regarding many aspects of the living animal, including feeding, movement, and social behaviour. As one of the first scientifically well-known dinosaurs, Iguanodon has occupied a small but notable place in the public's perception of dinosaurs, its artistic representation changing significantly in response to new interpretations of its remains.
Discovery and history
Gideon Mantell, Sir Richard Owen, and the discovery of dinosaurs
The discovery of Iguanodon has long been accompanied by a popular legend. The story goes that Gideon Mantell's wife, Mary Ann, discovered the first teeth of an Iguanodon in the strata of Tilgate Forest in Whitemans Green, Cuckfield, Sussex, England, in 1822 while her husband was visiting a patient. However, there is no evidence that Mantell took his wife with him while seeing patients. Furthermore, he admitted in 1851 that he himself had found the teeth, although he had previously stated in 1827 and 1833 that Mrs. Mantell had indeed found the first of the teeth later named Iguanodon. Other later authors agree that the story is not certainly false. It is known from his notebooks that Mantell first acquired large fossil bones from the quarry at Whitemans Green in 1820. Because also theropod teeth were found, thus belonging to carnivores, he at first interpreted these bones, which he tried to combine into a partial skeleton, as those of a giant crocodile. In 1821 Mantell mentioned the find of herbivorous teeth and began to consider the possibility that a large herbivorous reptile was present in the strata. However, in his 1822 publication Fossils of the South Downs he as yet did not dare to suggest a connection between the teeth and his very incomplete skeleton, presuming that his finds presented two large forms, one carnivorous ("an animal of the Lizard Tribe of enormous magnitude"), the other herbivorous.
In May 1822 he first presented the herbivorous teeth to the Royal Society of London but the members, among them William Buckland, dismissed them as fish teeth or the incisors of a rhinoceros from a Tertiary stratum. On 23 June 1823 Charles Lyell showed some to Georges Cuvier, during a soiree in Paris, but the famous French naturalist at once dismissed them as those of a rhinoceros. Though the very next day Cuvier retracted, Lyell reported only the dismissal to Mantell, who became rather diffident about the issue. In 1824 Buckland described Megalosaurus and was on that occasion invited to visit Mantell's collection. Seeing the bones on 6 March he agreed that these were of some giant saurian—though still denying it was a herbivore. Emboldened nevertheless, Mantell again sent some teeth to Cuvier, who answered on 22 June 1824 that he had determined that they were reptilian and quite possibly belonged to a giant herbivore. In a new edition that year of his Recherches sur les Ossemens Fossiles Cuvier admitted his earlier mistake, leading to an immediate acceptance of Mantell, and his new saurian, in scientific circles. Mantell tried to corroborate his theory further by finding a modern-day parallel among extant reptiles. In September 1824 he visited the Royal College of Surgeons but at first failed to find comparable teeth. However, assistant-curator Samuel Stutchbury recognised that they resembled those of an iguana he had recently prepared, albeit twenty times longer.
In recognition of the resemblance of the teeth to those of the iguana, Mantell decided to name his new animal Iguanodon or 'iguana-tooth', from iguana and the Greek word ὀδών (odon, odontos or 'tooth'). Based on isometric scaling, he estimated that the creature might have been up to long, more than the length of Megalosaurus. His initial idea for a name was Iguana-saurus ('Iguana lizard'), but his friend William Daniel Conybeare suggested that that name was more applicable to the iguana itself, so a better name would be Iguanoides ('Iguana-like') or Iguanodon. He neglected to add a specific name to form a proper binomial, but one was supplied in 1829 by Friedrich Holl: I. anglicum, which was later emended to I. anglicus.
Mantell sent a letter detailing his discovery to the local Portsmouth Philosophical Society in December 1824, several weeks after settling on a name for the fossil creature. The letter was read to members of the Society at a meeting on 17 December, and a report was published in the Hampshire Telegraph the following Monday, 20 December, which announced the name, misspelled as "Iguanadon". Mantell formally published his findings on 10 February 1825, when he presented a paper on the remains to the Royal Society of London.
A more complete specimen of a similar animal was discovered in a quarry in Maidstone, Kent, in 1834 (lower Lower Greensand Formation), which Mantell soon acquired. He was led to identify it as an Iguanodon based on its distinctive teeth. The Maidstone slab (NHMUK PV OR 3791) was used in the first skeletal reconstructions and artistic renderings of Iguanodon, but due to its incompleteness, Mantell made some mistakes, the most famous of which was the placement of what he thought was a horn on the nose. The discovery of much better specimens in later years revealed that the horn was actually a modified thumb. Still encased in rock, the Maidstone skeleton is currently displayed at the Natural History Museum in London. The borough of Maidstone commemorated this find by adding an Iguanodon as a supporter to their coat of arms in 1949. This specimen has become linked with the name I. mantelli, a species named in 1832 by Christian Erich Hermann von Meyer in place of I. anglicus, but it actually comes from a different formation than the original I. mantelli/I. anglicus material. The Maidstone specimen, also known as Gideon Mantell's "Mantel-piece", and formally labelled NHMUK 3741 was subsequently excluded from Iguanodon. It is classified as cf. Mantellisaurus by McDonald (2012); as cf. Mantellisaurus atherfieldensis by Norman (2012); and made the holotype of a separate species Mantellodon carpenteri by Paul (2012), but this is considered dubious and it is generally considered a specimen of Mantellisaurus
At the same time, tension began to build between Mantell and Richard Owen, an ambitious scientist with much better funding and society connections in the turbulent worlds of Reform Act-era British politics and science. Owen, a firm creationist, opposed the early versions of evolutionary science ("transmutationism") then being debated and used what he would soon coin as dinosaurs as a weapon in this conflict. With the paper describing Dinosauria, he scaled down dinosaurs from lengths of over , determined that they were not simply giant lizards, and put forward that they were advanced and mammal-like, characteristics given to them by God; according to the understanding of the time, they could not have been "transmuted" from reptiles to mammal-like creatures.
In 1849, a few years before his death in 1852, Mantell realised that iguanodonts were not heavy, pachyderm-like animals, as Owen was putting forward, but had slender forelimbs. However, since his passing left him unable to participate in the creation of the Crystal Palace dinosaur sculptures, Owen's vision of the dinosaurs became that seen by the public for decades. With Benjamin Waterhouse Hawkins, he had nearly two dozen lifesize sculptures of various prehistoric animals built out of concrete sculpted over a steel and brick framework; two iguanodonts (based on the Maidstone specimen), one standing and one resting on its belly, were included. Before the sculpture of the standing iguanodont was completed, he held a banquet for twenty inside it.
Bernissart mine discoveries and Dollo's new reconstruction
The largest find of Iguanodon remains to that date occurred on 28 February 1878 in a coal mine at Bernissart in Belgium, at a depth of , when two mineworkers, Jules Créteur and Alphonse Blanchard, accidentally hit on a skeleton that they initially took for petrified wood. With the encouragement of Alphonse Briart, supervisor of mines at nearby Morlanwelz, Louis de Pauw on 15 May 1878 started to excavate the skeletons and in 1882 Louis Dollo reconstructed them. At least 38 Iguanodon individuals were uncovered, most of which were adults. In 1882, the holotype specimen of I. bernissartensis became one of the first ever dinosaur skeletons mounted for display. It was put together in a chapel at the Palace of Charles of Lorraine using a series of adjustable ropes attached to scaffolding so that a lifelike pose could be achieved during the mounting process. This specimen, along with several others, first opened for public viewing in an inner courtyard of the palace in July 1883. In 1891 they were moved to the Royal Museum of Natural History, where they are still on display; nine are displayed as standing mounts, and nineteen more are still in the Museum's basement. The exhibit makes an impressive display in the Royal Belgian Institute of Natural Sciences, in Brussels. A replica of one of these is on display at the Oxford University Museum of Natural History and at the Sedgwick Museum in Cambridge. Most of the remains were referred to a new species, I. bernissartensis, a larger and much more robust animal than the English remains had yet revealed. One specimen, IRSNB 1551, was at first referred to the nebulous, gracile I. mantelli, but is currently referred to Mantellisaurus atherfieldensis. The skeletons were some of the first complete dinosaur skeletons known. Found with the dinosaur skeletons were the remains of plants, fish, and other reptiles, including the crocodyliform Bernissartia.
The science of conserving fossil remains was in its infancy, and new techniques had to be improvised to deal with what soon became known as "pyrite disease". Crystalline pyrite in the bones was being oxidized to iron sulphate, accompanied by an increase in volume that caused the remains to crack and crumble. When in the ground, the bones were isolated by anoxic moist clay that prevented this from happening, but when removed into the drier open air, the natural chemical conversion began to occur. To limit this effect, De Pauw immediately, in the mine-gallery, re-covered the dug-out fossils with wet clay, sealing them with paper and plaster reinforced by iron rings, forming in total about six hundred transportable blocks with a combined weight of a hundred and thirty tons. In Brussels after opening the plaster he impregnated the bones with boiling gelatine mixed with oil of cloves as a preservative. Removing most of the visible pyrite he then hardened them with hide glue, finishing with a final layer of tin foil. Damage was repaired with papier-mâché. This treatment had the unintended effect of sealing in moisture and extending the period of damage. In 1932 museum director Victor van Straelen decided that the specimens had to be completely restored again to safeguard their preservation. From December 1935 to August 1936 the staff at the museum in Brussels treated the problem with a combination of alcohol, arsenic, and 390 kilograms of shellac. This combination was intended to simultaneously penetrate the fossils (with alcohol), prevent the development of mold (with arsenic), and harden them (with shellac). The fossils entered a third round of conservation from 2003 until May 2007, when the shellac, hide glue and gelatine were removed and impregnated with polyvinyl acetate and cyanoacrylate and epoxy glues. Modern treatments of this problem typically involve either monitoring the humidity of fossil storage, or, for fresh specimens, preparing a special coating of polyethylene glycol that is then heated in a vacuum pump, so that moisture is immediately removed and pore spaces are infiltrated with polyethylene glycol to seal and strengthen the fossil.
Dollo's specimens allowed him to show that Owen's prehistoric pachyderms were not correct for Iguanodon. He instead modelled the skeletal mounts after the cassowary and wallaby, and put the spike that had been on the nose firmly on the thumb. His reconstruction would prevail for a long period of time, but would later be discounted.
Excavations at the quarry were stopped in 1881, although it was not exhausted of fossils, as recent drilling operations have shown. During World War I, when the town was occupied by German forces, preparations were made to reopen the mine for palaeontology, and Otto Jaekel was sent from Berlin to supervise. Just as the first fossiliferous layer was about to be uncovered, however, the German army surrendered and had to withdraw. Further attempts to reopen the mine were hindered by financial problems and were stopped altogether in 1921 when the mine flooded.
Turn of the century and the Dinosaur Renaissance
Research on Iguanodon decreased during the early part of the 20th century as World Wars and the Great Depression enveloped Europe. A new species that would become the subject of much study and taxonomic controversy, I. atherfieldensis, was named in 1925 by R. W. Hooley, for a specimen collected at Atherfield Point on the Isle of Wight.
Iguanodon was not part of the initial work of the dinosaur renaissance that began with the description of Deinonychus in 1969, but it was not neglected for long. David B. Weishampel's work on ornithopod feeding mechanisms provided a better understanding of how it fed, and David B. Norman's work on numerous aspects of the genus has made it one of the best-known dinosaurs. In addition, a further find of numerous disarticulated Iguanodon bones in Nehden, Nordrhein-Westphalen, Germany, has provided evidence for gregariousness in this genus, as the animals in this areally restricted find appear to have been killed by flash floods. At least 15 individuals, from long, have been found here, most of the individuals belong to the related Mantellisaurus (described as I. atherfieldensis, at that time believed to be another species of Iguanodon). but some are of I. bernissartensis.
One major revision to Iguanodon brought by the Renaissance would be another re-thinking of how to reconstruct the animal. A major flaw with Dollo's reconstruction was the bend he introduced into the tail. This organ was more or less straight, as shown by the skeletons he was excavating, and the presence of ossified tendons. In fact, to get the bend in the tail for a more wallaby or kangaroo-like posture, the tail would have had to be broken. With its correct, straight tail and back, the animal would have walked with its body held horizontal to the ground, arms in place to support the body if needed.
21st century research and the splitting of the genus
In the 21st century, Iguanodon material has been used in the search for dinosaur biomolecules. In research by Graham Embery et al., Iguanodon bones were processed to look for remnant proteins. In this research, identifiable remains of typical bone proteins, such as phosphoproteins and proteoglycans, were found in a rib. In 2007, Gregory S. Paul split I. atherfieldensis into a new, separate genus, Mantellisaurus which has been generally accepted. In 2009 fragmentary iguanodontid material was described from upper Barremian Paris Basin deposits in Auxerre, Burgundy. While not definitively diagnosable to the genus/species level, the specimen shares "obvious morphological and dimensional affinities" with I. bernissartensis.
In 2010, David Norman split the Valanginian species I. dawsoni and I. fittoni into Barilium and Hypselospinus respectively. After Norman 2010, over half a dozen new genera were named off English "Iguanodon" material. Carpenter and Ishida in 2010 named Proplanicoxa, Torilion and Sellacoxa while Gregory S. Paul in 2012 named Darwinsaurus, Huxleysaurus and Mantellodon and Macdonald et al. in 2012 named Kukufeldia. These species named after Norman 2010 are not considered valid and are considered various junior synonyms of Mantellisaurus, Barilium and Hypselospinus.
In 2011, a new genus Delapparentia was named for a specimen in Spain that was originally thought to belong to I. bernissartensis. The previous identification was subsequently reaffirmed in a new analysis of individual variation in the Belgian specimens, finding that the Delapparentia specimen was within the range of I. bernissartensis. In 2015 a new species of Iguanodon, I. galvensis, was named based on material including 13 juvenile (perinate) individuals found in the Camarillas Formation near Galve, Spain. In 2017 a new study was done of I. galvensis, with further evidence of distinctiveness from I. bernissartensis including several new autapomorphies. It was also found that the Delapparentia holotype (which is also from the Camarillas Formation) was not distinguishable from either I. bernissartensis or I. galvensis.
Description
Iguanodon were bulky herbivores that could shift from bipedality to quadrupedality. The only well-supported species, I. bernissartensis, is estimated to have measured about long as an adult, with some specimens possibly as long as , although this is likely an overestimate, given that the maximum body length of I. bernissartensis is reported to be . Although Gregory S. Paul suggested a body mass of on average, constructing a 3D mathematical model and employing allometry-based estimate suggests an I. bernissartensis close to long (smaller than average) weighs close to in body mass. Specimens of relatively large individuals have been reported in the 2020s: a specimen referred to as I. cf. galvensis was measured up to in length, while a new specimen of I. bernissartensis from the upper Barremian of the Iberian Peninsula was measured up to in length. Such large individuals would have weighed approximately .The arms of I. bernissartensis were long (up to 75% the length of the legs) and robust, with rather inflexible hands built so that the three central fingers could bear weight. The thumbs were conical spikes that stuck out away from the three main digits. In early restorations, the spike was placed on the animal's nose. Later fossils revealed the true nature of the thumb spikes, although their exact function is still debated. They could have been used for defense, or for foraging for food. The little finger was elongated and dextrous, and could have been used to manipulate objects. The phalangeal formula is 2-3-3-2-4, meaning that the innermost finger (phalange) has two bones, the next has three, etc. The legs were powerful, but not built for running, and each foot had three toes. The backbone and tail were supported and stiffened by ossified tendons, which were tendons that turned to bone during life (these rod-like bones are usually omitted from skeletal mounts and drawings).
These animals had large, tall but narrow skulls, with toothless beaks probably covered with keratin, and teeth like those of iguanas, as the name suggests, but much larger and more closely packed. Unlike hadrosaurids, which had columns of replacement teeth, Iguanodon only had one replacement tooth at a time for each position. The upper jaw held up to 29 teeth per side, with none at the front of the jaw, and the lower jaw 25; the numbers differ because teeth in the lower jaw are broader than those in the upper. Because the tooth rows are deeply inset from the outside of the jaws, and because of other anatomical details, it is believed that, as with most other ornithischians, Iguanodon had some sort of cheek-like structure, muscular or non-muscular, to retain food in the mouth.
Classification and evolution
Iguanodon gives its name to the unranked clade Iguanodontia, a very populous group of ornithopods with many species known from the Middle Jurassic to the Late Cretaceous. Aside from Iguanodon, the best-known members of the clade include Dryosaurus, Camptosaurus, Ouranosaurus, and the duck-bills, or hadrosaurs. In older sources, Iguanodontidae was shown as a distinct family. This family traditionally has been something of a wastebasket taxon, including ornithopods that were neither hypsilophodontids or hadrosaurids. In practice, animals like Callovosaurus, Camptosaurus, Craspedodon, Kangnasaurus, Mochlodon, Muttaburrasaurus, Ouranosaurus, and Probactrosaurus were usually assigned to this family.
With the advent of cladistic analyses, Iguanodontidae as traditionally construed was shown to be paraphyletic, and these animals are recognised to fall at different points in relation to hadrosaurs on a cladogram, instead of in a single distinct clade. Essentially, the modern concept of Iguanodontidae currently includes only Iguanodon. Groups like Iguanodontoidea are still used as unranked clades in the scientific literature, though many traditional iguanodontids are now included in the superfamily Hadrosauroidea. Iguanodon lies between Camptosaurus and Ouranosaurus in cladograms, and is probably descended from a camptosaur-like animal. At one point, Jack Horner suggested, based mostly on skull features, that hadrosaurids actually formed two more distantly related groups, with Iguanodon on the line to the flat-headed hadrosaurines, and Ouranosaurus on the line to the crested lambeosaurines, but his proposal has been rejected.
The cladogram below follows an analysis by Andrew McDonald, 2012.
Species
Because Iguanodon is one of the first dinosaur genera to have been named, numerous species have been assigned to it. While never becoming the wastebasket taxon several other early genera of dinosaurs (such as Megalosaurus) became, Iguanodon has had a complicated history, and its taxonomy continues to undergo revisions. Although Gregory Paul recommended restricting I. bernissartensis to the famous sample from Bernissart, ornithopod workers like Norman and McDonald have disagreed with Paul's recommendations, except exercising caution when accepting records of Iguanodon from France and Spain as valid.
I. anglicus was the original type species, but the lectotype was based on a single tooth and only partial remains of the species have been recovered since. In March 2000, the International Commission on Zoological Nomenclature changed the type species to the much better known I. bernissartensis, with the new holotype being IRSNB 1534. The original Iguanodon tooth is held at Te Papa Tongarewa, the national museum of New Zealand in Wellington, although it is not on display. The fossil arrived in New Zealand following the move of Gideon Mantell's son Walter there; after the elder Mantell's death, his fossils went to Walter.
Species currently accepted as valid
Only two species assigned to Iguanodon are still considered to be valid.
I. bernissartensis, described by George Albert Boulenger in 1881, is the type species for the genus. This species is best known for the many skeletons discovered in the Sainte-Barbe Clays Formation at Bernissart, but is also known from remains across Europe.
Delapparentia turolensis, named in 2011 based on a specimen previously assigned to Iguanodon bernissartensis, was argued to be distinct from the latter based on the relative height of its neural spines. However, a 2017 study noted that this is easily within the range of individual variation, and that the difference may also arise from D. turolensis being an adult older than other specimens of I. bernissartensis.
I. seelyi (also incorrectly spelled I. seeleyi), described by John Hulke in 1882, has also been synonymised with Iguanodon bernissartensis, though this is not universally accepted. It was discovered in Brook, on the Isle of Wight, and named after Charles Seely MP, Liberal politician and philanthropist, on whose estate it was found.
David Norman has suggested that I. bernissartensis includes the dubious Mongolian I. orientalis (see also below), but this has not been followed by other researchers.
I. galvensis, described in 2015, is based on adult and juvenile remains found in Barremian-age deposits in Teruel, Spain.
Reassigned species of Iguanodon
I. albinus (or Albisaurus scutifer), described by Czech palaeontologist Antonin Fritsch in 1893, is a dubious nondinosaurian reptile now known as Albisaurus albinus.
I. atherfieldensis, described by R.W. Hooley in 1925, was smaller and less robust than I. bernissartensis, with longer neural spines. It was renamed Mantellisaurus atherfieldensis in 2007. The Bernissart specimen RBINS 1551 was described as Dollodon bampingi in 2008, but McDonald and Norman returned Dollodon to synonymy with Mantellisaurus.
I. dawsoni, described by Lydekker in 1888, is known from two partial skeletons found in East Sussex, England, from the middle Valanginian-age Lower Cretaceous Wadhurst Clay. It is now the type species of Barilium.
I. exogyrarum was described by Fritsch in 1878. It is a nomen dubium based on very poor material and was renamed Ponerosteus in 2000.
I. fittoni was described by Lydekker in 1889. Like I. dawsoni, this species was described from the Wadhurst Clay of East Sussex. It is now the type species of Hypselospinus.
I. hilli, coined by Edwin Tully Newton in 1892 for a tooth from the early Cenomanian Upper Cretaceous Lower Chalk of Hertfordshire, has been considered an early hadrosaurid of some sort. However, recent work places it as indeterminate beyond Hadrosauroidea outside Hadrosauridae.
I. hoggi (also spelled I. boggii or hoggii), named by Owen for a lower jaw from the Tithonian–Berriasian-age Upper Jurassic–Lower Cretaceous Purbeck Beds of Dorset in 1874, has been reassigned to its own genus, Owenodon.
I. hollingtoniensis (also spelled I. hollingtonensis), described by Lydekker in 1889, has variously been considered a synonym of Hypselospinus fittoni or a distinct species assigned to the genus Huxleysaurus. A specimen from the Valanginian Wadhurst Clay Formation, variously assigned to I. hollingtoniensis and I. mantelli over the years, has an unusual combination of hadrosaurid-like lower jaw and very robust forelimb; Norman (2010) assigned this specimen to the species Hypselospinus fittoni, while Paul (2012) made it the holotype of a separate species Darwinsaurus evolutionis.
I. lakotaensis was described by David B. Weishampel and Philip R. Bjork in 1989. The only well-accepted North American species of Iguanodon, I. lakotaensis was described from a partial skull from the Barremian-age Lower Cretaceous Lakota Formation of South Dakota. Its assignment has been controversial. Some researchers suggest that it was more basal than I. bernissartensis, and related to Theiophytalia, but David Norman has suggested that it was a synonym of I. bernissartensis. Gregory S. Paul has since given the species its own genus, Dakotadon.
I. mantelli described by Christian Erich Hermann von Meyer in 1832, was based on the same material as I. anglicus and is an objective junior synonym of the latter. Several taxa, including the holotype of Mantellisaurus and Mantellodon, but also the dubious hadrosauroid Trachodon cantabrigiensis the hypsilophodont Hypsilophodon, and Valdosaurus, were previously mis-assigned to I. mantelli.
"I. mongolensis" is a nomen nudum from a photo caption in a book by Whitfield in 1992 of remains that would later be named Altirhinus.
I. orientalis, described by A. K. Rozhdestvensky in 1952, was based on poor material, but a skull with a distinctive arched snout that had been assigned to it was renamed Altirhinus kurzanovi in 1998. At the same time, I. orientalis was considered to be a nomen dubium because it cannot be compared to I. bernissartensis.
I. phillipsi was described by Harry Seeley in 1869, but he later reassigned it to Priodontognathus.
I. praecursor (also spelled I. precursor), described by E. Sauvage in 1876 from teeth from an unnamed Kimmeridgian (Late Jurassic) formation in Pas-de-Calais, France, is actually a sauropod, sometimes assigned to Neosodon, although the two come from different formations.
I. prestwichii (also spelled I. prestwichi), described by John Hulke in 1880, has been reassigned to Camptosaurus prestwichii or to its own genus Cumnoria.
I. suessii, described by Emanuel Bunzel in 1871, has been reassigned to Mochlodon suessi.
Species reassigned to Iguanodon
I. foxii (also spelled I. foxi) was originally described by Thomas Henry Huxley in 1869 as the type species of Hypsilophodon; Owen (1873 or 1874) reassigned it to Iguanodon, but his assignment was soon overturned.
I. gracilis, named by Lydekker in 1888 as the type species of Sphenospondylus and assigned to Iguanodon in 1969 by Rodney Steel, has been suggested to be a synonym of Mantellisaurus atherfieldensis, but is considered dubious nowadays.
I. major, a species named by Justin Delair in 1966, based on vertebrae from the Isle of Wight and Sussex originally described by Owen in 1842 as a species of Streptospondylus, S. major, is a nomen dubium.
I. valdensis, a renaming of Vectisaurus valdensis by Ernst van den Broeck in 1900. Originally named by Hulke as a distinct genus in 1879 based on vertebral and pelvic remains, it was from the Barremian stage of the Isle of Wight. It was considered a juvenile specimen of Mantellisaurus atherfieldensis, or an undetermined species of Mantellisaurus, but is indeterminate beyond Iguanodontia.
The nomen nudum "Proiguanodon" (van den Broeck, 1900) also belongs here.
Dubious species
I. anglicus, described by Friedrich Holl in 1829, is the original type species of Iguanodon, but, as discussed above, was replaced by I. bernissartensis. In the past, it has been spelled as I. angelicus (Lessem and Glut, 1993) and I. anglicum (Holl, 1829 emend. Bronn, 1850). It is possible teeth ascribed to this species belong to the genus now called Barilium. The name Therosaurus (Fitzinger, 1840), is a junior objective synonym, a later name for the material of I. anglicus.
I. ottingeri, described by Peter Galton and James A. Jensen in 1979, is a nomen dubium based on teeth from the possibly Aptian-age lower Cedar Mountain Formation of Utah.
Palaeobiology
Feeding
One of the first details noted about Iguanodon was that it had the teeth of a herbivorous reptile, although there has not always been consensus on how it ate. As Mantell noted, the remains he was working with were unlike any modern reptile, especially in the toothless, scoop-shaped form of the lower jaw symphysis, which he found best compared to that of the two-toed sloth and the extinct ground sloth Mylodon. He also suggested that Iguanodon had a prehensile tongue which could be used to gather food, like a giraffe. More complete remains have shown this to be an error; for example, the hyoid bones that supported the tongue are heavily built, implying a muscular, non-prehensile tongue used for moving food around in the mouth. The giraffe-tongue idea has also been incorrectly attributed to Dollo via a broken lower jaw.
The skull was structured in such a way that as it closed, the bones holding the teeth in the upper jaw would bow out. This would cause the lower surfaces of the upper jaw teeth to rub against the upper surface of the lower jaw's teeth, grinding anything caught in between and providing an action that is the rough equivalent of mammalian chewing. Because the teeth were always replaced, the animal could have used this mechanism throughout its life, and could eat tough plant material. Additionally, the front ends of the animal's jaws were toothless and tipped with bony nodes, both upper and lower, providing a rough margin that was likely covered and lengthened by a keratinous material to form a cropping beak for biting off twigs and shoots. Its food gathering would have been aided by its flexible little finger, which could have been used to manipulate objects, unlike the other fingers.
Exactly what Iguanodon ate with its well-developed jaws is not known. The size of the larger species, such as I. bernissartensis, would have allowed them access to food from ground level to tree foliage at high. A diet of horsetails, cycads, and conifers was suggested by David Norman, although iguanodonts in general have been tied to the advance of angiosperm plants in the Cretaceous due to the dinosaurs' inferred low-browsing habits. Angiosperm growth, according to this hypothesis, would have been encouraged by iguanodont feeding because gymnosperms would be removed, allowing more space for the weed-like early angiosperms to grow. The evidence is not conclusive, though. Whatever its exact diet, due to its size and abundance, Iguanodon is regarded as a dominant medium to large herbivore for its ecological communities. In England, this included the small predator Aristosuchus, larger predators Eotyrannus, Baryonyx, and Neovenator, low-feeding herbivores Hypsilophodon and Valdosaurus, fellow "iguanodontid" Mantellisaurus, the armoured herbivore Polacanthus, and sauropods like Pelorosaurus.
Posture and movement
Early fossil remains were fragmentary, which led to much speculation on the posture and nature of Iguanodon. Iguanodon was initially portrayed as a quadrupedal horn-nosed beast. However, as more bones were discovered, Mantell observed that the forelimbs were much smaller than the hindlimbs. His rival Owen was of the opinion it was a stumpy creature with four pillar-like legs. The job of overseeing the first lifesize reconstruction of dinosaurs was initially offered to Mantell, who declined due to poor health, and Owen's vision subsequently formed the basis on which the sculptures took shape. Its bipedal nature was revealed with the discovery of the Bernissart skeletons. However, it was depicted in an upright posture, with the tail dragging along the ground, acting as the third leg of a tripod.
During his re-examination of Iguanodon, David Norman was able to show that this posture was unlikely, because the long tail was stiffened with ossified tendons. To get the tripodal pose, the tail would literally have to be broken. Putting the animal in a horizontal posture makes many aspects of the arms and pectoral girdle more understandable. For example, the hand is relatively immobile, with the three central fingers grouped together, bearing hoof-like phalanges, and able to hyperextend. This would have allowed them to bear weight. The wrist is also relatively immobile, and the arms and shoulder bones robust. These features all suggest that the animal spent time on all fours.
Furthermore, it appears that Iguanodon became more quadrupedal as it got older and heavier; juvenile I. bernissartensis have shorter arms than adults (60% of hindlimb length versus 70% for adults). When walking as a quadruped, the animal's hands would have been held so that the palms faced each other, as shown by iguanodontian trackways and the anatomy of this genus's arms and hands. The three-toed pes (foot) of Iguanodon was relatively long, and when walking, both the hand and the foot would have been used in a digitigrade fashion (walking on the fingers and toes). The maximum speed of Iguanodon has been estimated at , which would have been as a biped; it would not have been able to gallop as a quadruped.
Large three-toed footprints are known in Early Cretaceous rocks of England, particularly Wealden beds on the Isle of Wight, and these trace fossils were originally difficult to interpret. Some authors associated them with dinosaurs early on. In 1846, E. Tagert went so far as to assign them to an ichnogenus he named Iguanodon, and Samuel Beckles noted in 1854 that they looked like bird tracks, but might have come from dinosaurs. The identity of the trackmakers was greatly clarified upon the discovery in 1857 of the hind leg of a young Iguanodon, with distinctly three-toed feet, showing that such dinosaurs could have made the tracks. Despite the lack of direct evidence, these tracks are often attributed to Iguanodon. A trackway in England shows what may be an Iguanodon moving on all fours, but the foot prints are poor, making a direct connection difficult. Tracks assigned to the ichnogenus Iguanodon are known from locations including places in Europe where the body fossil Iguanodon is known, to Spitsbergen, Svalbard, Norway.
Thumb spike
The thumb spike is one of the best-known features of Iguanodon. Although it was originally placed on the animal's nose by Mantell, the complete Bernissart specimens allowed Dollo to place it correctly on the hand, as a modified thumb. (This would not be the last time a dinosaur's modified thumb claw would be misinterpreted; Noasaurus, Baryonyx, and Megaraptor are examples since the 1980s where an enlarged thumb claw was first put on the foot, as in dromaeosaurids.)
This thumb is typically interpreted as a close-quarter stiletto-like weapon against predators, although it could also have been used to break into seeds and fruits, or against other Iguanodon. One author has suggested that the spike was attached to a venom gland, but this has not been accepted, as the spike was not hollow, nor were there any grooves on the spike for conducting venom.
Possible social behaviour
Although sometimes interpreted as the result of a single catastrophe, the Bernissart finds instead are now interpreted as recording multiple events. According to this interpretation, at least three occasions of mortality are recorded, and though numerous individuals would have died in a geologically short time span (?10–100 years), this does not necessarily mean these Iguanodon were herding animals.
An argument against herding is that juvenile remains are very uncommon at this site, unlike modern cases with herd mortality. They more likely were the periodic victims of flash floods whose carcasses accumulated in a lake or marshy setting. The Nehden find, however, with its greater span of individual ages, more even mix of Dollodon or Mantellisaurus to Iguanodon bernissartensis, and confined geographic nature, may record mortality of herding animals migrating through rivers.
There is no evidence that Iguanodon was sexually dimorphic (with one sex appreciably different from the other). At one time, it was suggested that the Bernissart I. "mantelli", or I. atherfieldensis (Dollodon and Mantellisaurus, respectively) represented a sex, possibly female, of the larger and more robust, possibly male, I. bernissartensis. However, this is not supported today. A 2017 analysis showed that I. bernissartensis does exhibit a large level of individual variation in both its limbs (scapula, humerus, thumb claw, ilium, ischium, femur, tibia) and spinal column (axis, sacrum, tail vertebrae). Additionally, this analysis found that individuals of I. bernissartensis generally seemed to fall into two categories based on whether their tail vertebrae bore a furrow on the bottom, and whether their thumb claws were large or small.
Paleopathology
Evidence of a fractured hip bone was found in a specimen of Iguanodon, which had an injury to its ischium. Two other individuals were observed with signs of osteoarthritis as evidenced by bone overgrowths in their anklebones which are called osteophytes.
In popular culture
Since its description in 1825, Iguanodon has been a feature of worldwide popular culture. Two lifesize reconstructions of Mantellodon (considered Iguanodon at the time) built at the Crystal Palace in London in 1852 greatly contributed to the popularity of the genus. Their thumb spikes were mistaken for horns, and they were depicted as elephant-like quadrupeds, yet this was the first time an attempt was made at constructing full-size dinosaur models. In 1910 Heinrich Harder portrayed a group of Iguanodon in Tiere der Urwelt, a classic German collecting card game about extinct and prehistoric animals.
Several motion pictures have featured Iguanodon. In the 2000 Disney animated film Dinosaur, an Iguanodon named Aladar served as the protagonist with three other iguanodonts as other main and minor characters are Neera, Kron and Bruton. A loosely related ride of the same name at Disney's Animal Kingdom is based around bringing an Iguanodon back to the present. Iguanodon is one of the three dinosaur genera that inspired Godzilla; the other two were Tyrannosaurus rex and Stegosaurus. Iguanodon has also made appearances in some of the many The Land Before Time films, as well as episodes of the television series.
Aside from appearances in movies, Iguanodon has also been featured on the television documentary miniseries Walking with Dinosaurs (1999) produced by the BBC (along with then-undescribed Dakotadon lakotaensis) and played a starring role in Sir Arthur Conan Doyle's book The Lost World as well as featuring in the 2015 documentary Dinosaur Britain. It also was present in Bob Bakker's Raptor Red (1995), as a Utahraptor prey item. A main belt asteroid, , has been named 9941 Iguanodon in honour of the genus.
Because it is both one of the first dinosaurs described and one of the best-known dinosaurs, Iguanodon has been well-placed as a barometer of changing public and scientific perceptions on dinosaurs. Its reconstructions have gone through three stages: the elephantine quadrupedal horn-snouted reptile satisfied the Victorians, then a bipedal but still fundamentally reptilian animal using its tail to prop itself up dominated the early 20th century, but was slowly overturned during the 1960s by its current, more agile and dynamic representation, able to shift from two legs to all fours.
| Biology and health sciences | Dinosaurs and prehistoric reptiles | null |
228839 | https://en.wikipedia.org/wiki/Blazar | Blazar | A blazar is an active galactic nucleus (AGN) with a relativistic jet (a jet composed of ionized matter traveling at nearly the speed of light) directed very nearly towards an observer. Relativistic beaming of electromagnetic radiation from the jet makes blazars appear much brighter than they would be if the jet were pointed in a direction away from Earth. Blazars are powerful sources of emission across the electromagnetic spectrum and are observed to be sources of high-energy gamma ray photons. Blazars are highly variable sources, often undergoing rapid and dramatic fluctuations in brightness on short timescales (hours to days). Some blazar jets appear to exhibit superluminal motion, another consequence of material in the jet traveling toward the observer at nearly the speed of light.
The blazar category includes BL Lac objects and optically violently variable (OVV) quasars. The generally accepted theory is that BL Lac objects are intrinsically low-power radio galaxies while OVV quasars are intrinsically powerful radio-loud quasars. The name "blazar" was coined in 1978 by astronomer Edward Spiegel to denote the combination of these two classes.
In visible-wavelength images, most blazars appear compact and pointlike, but high-resolution images reveal that they are located at the centers of elliptical galaxies.
Blazars are important topics of research in astronomy and high-energy astrophysics. Blazar research includes investigation of the properties of accretion disks and jets, the central supermassive black holes and surrounding host galaxies, and the emission of high-energy photons, cosmic rays, and neutrinos.
In July 2018, the IceCube Neutrino Observatory team traced a neutrino that hit its Antarctica-based detector in September 2017 to its point of origin in a blazar 3.7 billion light-years away. This was the first time that a neutrino detector was used to locate an object in space.
Structure
Blazars, like all active galactic nuclei (AGN), are thought to be powered by material falling into a supermassive black hole in the core of the host galaxy. Gas, dust and the occasional star are captured and spiral into this central black hole, creating a hot accretion disk which generates enormous amounts of energy in the form of photons, electrons, positrons and other elementary particles. This region is relatively small, approximately 10−3 parsecs in size.
There is also a larger opaque toroid extending several parsecs from the black hole, containing a hot gas with embedded regions of higher density. These "clouds" can absorb and re-emit energy from regions closer to the black hole. On Earth, the clouds are detected as emission lines in the blazar spectrum.
Perpendicular to the accretion disk, a pair of relativistic jets carries highly energetic plasma away from the AGN. The jet is collimated by a combination of intense magnetic fields and powerful winds from the accretion disk and toroid. Inside the jet, high energy photons and particles interact with each other and the strong magnetic field. These relativistic jets can extend as far as many tens of kiloparsecs from the central black hole.
All of these regions can produce a variety of observed energy, mostly in the form of a nonthermal spectrum ranging from very low-frequency radio to extremely energetic gamma rays, with a high polarization (typically a few percent) at some frequencies. The nonthermal spectrum consists of synchrotron radiation in the radio to X-ray range, and inverse Compton emission in the X-ray to gamma-ray region. A thermal spectrum peaking in the ultraviolet region and faint optical emission lines are also present in OVV quasars, but faint or non-existent in BL Lac objects.
Relativistic beaming
The observed emission from a blazar is greatly enhanced by relativistic effects in the jet, a process called relativistic beaming. The bulk speed of the plasma that constitutes the jet can be in the range of 95%–99% of the speed of light, although individual particles move at higher speeds in various directions.
The relationship between the luminosity emitted in the rest frame of the jet and the luminosity observed from Earth depends on the characteristics of the jet. These include whether the luminosity arises from a shock front or a series of brighter blobs in the jet, as well as details of the magnetic fields within the jet and their interaction with the moving particles.
A simple model of beaming illustrates the basic relativistic effects connecting the luminosity in the rest frame of the jet, Se, and the luminosity observed on Earth, So: So is proportional to Se × D2, where D is the doppler factor.
When considered in much more detail, three relativistic effects are involved:
Relativistic aberration contributes a factor of D2. Aberration is a consequence of special relativity where directions which appear isotropic in the rest frame (in this case, the jet) appear pushed towards the direction of motion in the observer's frame (in this case, Earth).
Time dilation contributes a factor of D+1. This effect speeds up the apparent release of energy. If the jet emits a burst of energy every minute in its own rest frame, this release would be observed on Earth as much more frequent, perhaps every ten seconds.
Windowing can contribute a factor of D−1 and then works to decrease boosting. This happens for a steady flow because there are then D fewer elements of fluid within the observed window, as each element has been expanded by factor D. However, for a freely propagating blob of material, the radiation is boosted by the full D+3.
Example
Consider a jet with an angle to the line of sight θ = 5° and a speed of 99.9% of the speed of light. The luminosity observed from Earth is 70 times greater than the emitted luminosity. However, if θ is at the minimum value of 0° the jet will appear 600 times brighter from Earth.
Beaming away
Relativistic beaming also has another critical consequence. The jet which is not approaching Earth will appear dimmer because of the same relativistic effects. Therefore, two intrinsically identical jets will appear significantly asymmetric. In the example given above any jet where θ > 35° will be observed on Earth as less luminous than it would be from the rest frame of the jet.
A further consequence is that a population of intrinsically identical AGN scattered in space with random jet orientations will look like a very inhomogeneous population on Earth. The few objects where θ is small will have one very bright jet, while the rest will apparently have considerably weaker jets. Those where θ varies from 90° will appear to have asymmetric jets.
This is the essence behind the connection between blazars and radio galaxies. AGN which have jets oriented close to the line of sight with Earth can appear extremely different from other AGN even if they are intrinsically identical.
Discovery
Many of the brighter blazars were first identified, not as powerful distant galaxies, but as irregular variable stars in our own galaxy. These blazars, like genuine irregular variable stars, changed in brightness on periods of days or years, but with no pattern.
The early development of radio astronomy had shown that there are many bright radio sources in the sky. By the end of the 1950s, the resolution of radio telescopes was sufficient to identify specific radio sources with optical counterparts, leading to the discovery of quasars. Blazars were highly represented among these early quasars, and the first redshift was found for 3C 273, a highly variable quasar which is also a blazar.
In 1968, a similar connection was made between the "variable star" BL Lacertae and a powerful radio source VRO 42.22.01. BL Lacertae shows many of the characteristics of quasars, but the optical spectrum was devoid of the spectral lines used to determine redshift. Faint indications of an underlying galaxy—proof that BL Lacertae was not a star—were found in 1974.
The extragalactic nature of BL Lacertae was not a surprise. In 1972 a few variable optical and radio sources were grouped together and proposed as a new class of galaxy: BL Lacertae-type objects. This terminology was soon shortened to "BL Lacertae object", "BL Lac object" or simply "BL Lac". (The latter term can also mean the original individual blazar and not the entire class.)
, a few hundred BL Lac objects were known. One of the closest blazars is 2.5 billion light years away.
Current view
Blazars are thought to be active galactic nuclei, with relativistic jets oriented close to the line of sight with the observer.
The special jet orientation explains the general peculiar characteristics: high observed luminosity, very rapid variation, high polarization (compared to non-blazar quasars), and the apparent superluminal motions detected along the first few parsecs of the jets in most blazars.
A Unified Scheme or Unified Model has become generally accepted, where highly variable quasars are related to intrinsically powerful radio galaxies, and BL Lac objects are related to intrinsically weak radio galaxies. The distinction between these two connected populations explains the difference in emission line properties in blazars.
Other explanations for the relativistic jet/unified scheme approach which have been proposed include gravitational microlensing and coherent emission from the relativistic jet. Neither of these explains the overall properties of blazars. For example, microlensing is achromatic. That is, all parts of a spectrum would rise and fall together. This is not observed in blazars. However, it is possible that these processes, as well as more complex plasma physics, can account for specific observations or some details.
Examples of blazars include 3C 454.3, 3C 273, BL Lacertae, PKS 2155-304, Markarian 421, Markarian 501, 4C +71.07, PKS 0537-286 (QSO 0537-286) and S5 0014+81. Markarian 501 and S5 0014+81 are also called "TeV Blazars" for their high energy (teraelectron-volt range) gamma-ray emission.
In July 2018, a blazar called TXS 0506+056 was identified as source of high-energy neutrinos by the IceCube project.
| Physical sciences | Active galactic nucleus | null |
228845 | https://en.wikipedia.org/wiki/Connective%20tissue | Connective tissue | Connective tissue is one of the four primary types of animal tissue, a group of cells that are similar in structure, along with epithelial tissue, muscle tissue, and nervous tissue. It develops mostly from the mesenchyme, derived from the mesoderm, the middle embryonic germ layer. Connective tissue is found in between other tissues everywhere in the body, including the nervous system. The three meninges, membranes that envelop the brain and spinal cord, are composed of connective tissue. Most types of connective tissue consists of three main components: elastic and collagen fibers, ground substance, and cells. Blood, and lymph are classed as specialized fluid connective tissues that do not contain fiber. All are immersed in the body water. The cells of connective tissue include fibroblasts, adipocytes, macrophages, mast cells and leukocytes.
The term "connective tissue" (in German, ) was introduced in 1830 by Johannes Peter Müller. The tissue was already recognized as a distinct class in the 18th century.
Types
Connective tissue can be broadly classified into connective tissue proper, and special connective tissue. Connective tissue proper includes loose connective tissue, and dense connective tissue. Loose and dense connective tissue are distinguished by the ratio of ground substance to fibrous tissue. Loose connective tissue has much more ground substance and a relative lack of fibrous tissue, while the reverse is true of dense connective tissue.
Loose connective tissue
Loose connective tissue includes reticular connective tissue, and adipose tissue.
Dense connective tissue
Dense connective tissue also known as fibrous tissue is subdivided into dense regular and dense irregular connective tissue. Dense regular connective tissue, found in structures such as tendons and ligaments, is characterized by collagen fibers arranged in an orderly parallel fashion, giving it tensile strength in one direction. Dense irregular connective tissue provides strength in multiple directions by its dense bundles of fibers arranged in all directions.
Special connective tissue
Special connective tissue consists of cartilage, bone, blood and lymph. Other kinds of connective tissues include fibrous, elastic, and lymphoid connective tissues. Fibroareolar tissue is a mix of fibrous and areolar tissue. Fibromuscular tissue is made up of fibrous tissue and muscular tissue. New vascularised connective tissue that forms in the process of wound healing is termed granulation tissue. All of the special connective tissue types have been included as a subset of fascia in the fascial system, with blood and lymph classed as liquid fascia.
Bone and cartilage can be further classified as supportive connective tissue. Blood and lymph can also be categorized as fluid connective tissue, and liquid fascia.
Membranes
Membranes can be either of connective tissue or epithelial tissue. Connective tissue membranes include the meninges (the three membranes covering the brain and spinal cord) and synovial membranes that line joint cavities. Mucous membranes and serous membranes are epithelial with an underlying layer of loose connective tissue.
Fibrous types
Fiber types found in the extracellular matrix are collagen fibers, elastic fibers, and reticular fibers.
Ground substance is a clear, colorless, and viscous fluid containing glycosaminoglycans and proteoglycans allowing fixation of Collagen fibers in intercellular spaces. Examples of non-fibrous connective tissue include adipose tissue (fat) and blood. Adipose tissue gives "mechanical cushioning" to the body, among other functions. Although there is no dense collagen network in adipose tissue, groups of adipose cells are kept together by collagen fibers and collagen sheets in order to keep fat tissue under compression in place (for example, the sole of the foot). Both the ground substance and proteins (fibers) create the matrix for connective tissue.
Type I collagen is present in many forms of connective tissue, and makes up about 25% of the total protein content of the mammalian body.
Function
Connective tissue has a wide variety of functions that depend on the types of cells and the different classes of fibers involved. Loose and dense irregular connective tissue, formed mainly by fibroblasts and collagen fibers, have an important role in providing a medium for oxygen and nutrients to diffuse from capillaries to cells, and carbon dioxide and waste substances to diffuse from cells back into circulation. They also allow organs to resist stretching and tearing forces. Dense regular connective tissue, which forms organized structures, is a major functional component of tendons, ligaments and aponeuroses, and is also found in highly specialized organs such as the cornea. Elastic fibers, made from elastin and fibrillin, also provide resistance to stretch forces. They are found in the walls of large blood vessels and in certain ligaments, particularly in the ligamenta flava.
In hematopoietic and lymphatic tissues, reticular fibers made by reticular cells provide the stroma—or structural support—for the parenchyma (that is, the bulk of functional substance) of the organ.
Mesenchyme is a type of connective tissue found in developing organs of embryos that is capable of differentiation into all types of mature connective tissue. Another type of relatively undifferentiated connective tissue is the mucous connective tissue known as Wharton's jelly, found inside the umbilical cord. This tissue is no longer present after birth, leaving only scattered mesenchymal cells throughout the body.
Various types of specialized tissues and cells are classified under the spectrum of connective tissue, and are as diverse as brown and white adipose tissue, blood, cartilage and bone. Cells of the immune system—such as macrophages, mast cells, plasma cells, and eosinophils—are found scattered in loose connective tissue, providing the ground for starting inflammatory and immune responses upon the detection of antigens.
Clinical significance
There are many types of connective tissue disorders, such as:
Connective tissue neoplasms including sarcomas such as hemangiopericytoma and malignant peripheral nerve sheath tumor in nervous tissue.
Congenital diseases include Marfan syndrome and Ehlers-Danlos Syndrome.
Myxomatous degeneration – a pathological weakening of connective tissue.
Mixed connective tissue disease – a disease of the autoimmune system, also undifferentiated connective tissue disease.
Systemic lupus erythematosus (SLE) – a major autoimmune disease of connective tissue
Scurvy, caused by a deficiency of vitamin C which is necessary for the synthesis of collagen.
Fibromuscular dysplasia is a disease of the blood vessels that leads to an abnormal growth in the arterial wall.
| Biology and health sciences | Tissues | null |
228907 | https://en.wikipedia.org/wiki/Geoduck | Geoduck | The Pacific geoduck ( ; Panopea generosa) is a species of very large saltwater clam in the family Hiatellidae. The common name is derived from the Lushootseed name, .
The geoduck is native to the coastal waters of the eastern North Pacific Ocean from Alaska to Baja California. The shell of the clam ranges from to over in length, but the extremely long siphons make the clam itself much longer than this: the "shaft" or siphons alone can be in length. The geoduck is the largest burrowing clam in the world. It is also one of the longest-living animals of any type, with a typical lifespan of 140 years; the oldest has been recorded at 179 years old. The precise longevity of geoducks can be determined from annual rings deposited in the shell which can be assigned to calendar years of formation through crossdating. These annual rings also serve as an archive of past marine variability.
Etymology
The name Geoduck is derived from the Lushootseed name for the animal, . The etymology of is disputed. The lexical suffix means "many" in Lushootseed. The Oxford English Dictionary claims it is composed of a root word of unknown meaning and instead meaning "genitals" (referring to the shape of the clam), while other researchers claim it is a phrase meaning "dig deep".
It is sometimes known as a mud duck, king clam or, when translated literally from Chinese, an elephant-trunk clam ().
Between 1983 and 2010, the scientific name of this clam was confused with that of an extinct clam, Panopea abrupta (Conrad, 1849), in scientific literature.
Biology
Native to the west coast of Canada and the northwest coast of the United States (primarily Washington and British Columbia), these marine bivalve mollusks are the largest burrowing clams in the world, weighing in at an average of at maturity, but specimens weighing over and as much as in length are not unheard of.
A related species, Panopea zelandica, is found in New Zealand and has been harvested commercially since 1989. The largest quantities have come from Golden Bay in the South Island where were harvested in one year. There is a growing concern over the increase of parasites in the Puget Sound population of geoduck. Whether these microsporidium-like parasitic species were introduced by commercial farming is being studied by Sea Grant. Research to date does indicate their presence.
The oldest recorded specimen was 179 years old, but individuals usually live up to 140 years. A geoduck sucks water containing plankton down through its long siphon, filters this for food and ejects its refuse out through a separate hole in the siphon. Adult geoducks have few natural predators, which may also contribute to their longevity. In Alaska, sea otters and dogfish have proved capable of dislodging geoducks; starfish also attack and feed on the exposed geoduck siphon.
Geoducks are broadcast spawners. A female geoduck produces about 5 billion eggs in her century-long lifespan. However, due to a low rate of recruitment and a high rate of mortality for geoduck eggs, larvae, and post-settled juveniles, populations are slow to rebound. In the Puget Sound, studies indicate that the recovery time for a harvested tract is 39 years.
Biomass densities in Southeast Alaska are estimated by divers, then inflated by twenty percent to account for geoducks not visible at the time of survey. This estimate is used to predict the two percent allowed for commercial harvesting.
Industry
The world's first geoduck fishery was created in 1970, but demand for the half-forgotten clam was low at first due to its texture. , these clams sell in China for over US.
The geoduck's high market value has created an $80-million industry, with harvesting occurring in the US states of Alaska, Washington, and Oregon and the Canadian province of British Columbia. It is one of the most closely regulated fisheries in both countries. In Washington, Department of Natural Resources staff are on the water continually monitoring harvests to ensure revenues are received, and the same is true in Canada where the Underwater Harvesters' Association manages the Canadian Fishery in conjunction with Canada's Department of Fisheries and Oceans. The Washington State Department of Health tests water and flesh to assure clams are not filtering and holding pollutants, an ongoing problem. With the rise in price has come the inevitable problem with poaching, and with it the possibility some could be harvested from unsafe areas.
, advances in the testing system for contaminated clams have allowed geoduck harvesters to deliver live clams more consistently. The new testing system determines the viability of clams from tested beds before the harvesters fish the area. Previous methods tested clams after harvest. This advancement has meant that 90 percent of clams were delivered live to market in 2007. In 2001, only 10 percent were live. Because geoduck have a much higher market value live, an additional , this development has helped to stimulate the burgeoning industry.
The COVID-19 pandemic disrupted the geoduck industry. Given the near-shutdown of restaurants and seafood markets across the country, demand for live geoducks plummeted. Divers in Southeast Alaska who typically see prices of for live geoducks reported prices as low as , leading many to stop fishing temporarily.
Environmental impact
Geoduck farming grow-out and harvest practices are controversial, and have created conflicts with shoreline property owners, and concerns from nongovernmental organizations. However, the Environmental Defense Fund has found that bivalves (oysters, mussels, and clams) are beneficial to the marine environment. The water must be certifiably clean to plant geoducks commercially. Regulation was mandated in 2007. Studies have been funded to determine short- and long-term environmental and genetic impacts. In southern Puget Sound, the effect of geoduck farming on large mobile animals is ambiguous. A 2004 draft biological assessment, commissioned by three of the largest commercial shellfish companies in the Puget Sound region, identified no long-term effects of geoduck farming on threatened or endangered species.
Culinary uses
The large, meaty siphon is prized for its savory flavor and crunchy texture. Geoduck is regarded by some as an aphrodisiac because of its phallic shape. It is very popular in China, where it is considered a delicacy, mostly eaten cooked in a fondue-style Chinese hot pot. In Korean cuisine, geoducks are eaten raw with spicy chili sauce, sautéed, or in soups and stews. In Japan, geoduck is prepared as raw sashimi, dipped in soy sauce and wasabi. On Japanese menus in cheaper sushi restaurants, geoduck is sometimes substituted for Tresus keenae, a species of horse clam, and labeled or . It is considered to have a texture similar to an ark shell (known in Japanese as ). is sometimes translated into English as "giant clam", and it is distinguished from sushi, which is made from Tridacna gigas.
In popular culture
Evergreen State College in Olympia, Washington, has a geoduck as its mascot named Speedy.
Geoducks have also earned some internet infamy due to the phallic appearance of their siphons.
| Biology and health sciences | Bivalvia | Animals |
228964 | https://en.wikipedia.org/wiki/Taiwan%20High%20Speed%20Rail | Taiwan High Speed Rail | Taiwan High Speed Rail (THSR) is the high-speed railway network in Taiwan, which consists of a single line that runs approximately along the western coast of the island, from the capital Taipei in the north to the southern city of Kaohsiung. With construction and operations managed by a private company, Taiwan High Speed Rail Corporation (THSRC), which also operates the line, the total cost of the project was billion in 1998. The system's technology is based primarily on Japan's Shinkansen.
The railway opened for service on 5 January 2007, with trains running at a top speed of , currently running from Nangang to Zuoying in as little as 1 hour and 45 minutes, reaching almost 90% of Taiwan's population. Most intermediate stations on the line lie outside the cities served; however, a variety of transfer options, such as free shuttle buses, conventional rail, and metros have been constructed to facilitate transport connections.
Ridership initially fell short of forecasts, but grew from fewer than 40,000 passengers per day in the first few months of operation to over 129,000 passengers per day in June 2013. Daily passenger traffic reached 130,000 in 2014, well below the forecast of 240,000 daily passengers for 2008. The system carried its first 100 million passengers by August 2010 and over 200 million passengers had taken the system by December 2012, followed by 400 million by December 2016. THSR, or railways in general, is only located on the main island of Taiwan. Outer islands under the control of the ROC government including Penghu, Kinmen, and Matsu do not have railways.
In the initial years of operation, THSRC accumulated high debts due to high depreciation charges and interest, largely due to the financial structure set up for the private company. In 2009, THSRC negotiated with the government to change the method of depreciation from depending on concessions on rights to ridership. At the same time, the government also started to help refinance THSRC's loans to assist the company so it could remain operational and profitable. The government injected NT$30 billion as a financial bailout, boosting the government's stake to about 64% from about 37%. The government also extended the rail concession from 35 years to 70 years and terminated the company's build-operate-transfer business model.
History
Taiwan's rapid economic growth during the latter half of the twentieth century led to congestion of highways, conventional rail, and air traffic systems in the western transport corridor, which threatened to impede the region's development. The idea of a new high-speed rail line arose in the 1970s, and informal planning began in 1980. In 1987, the executive branch of Taiwan's government, the Executive Yuan, instructed the Ministry of Transportation to launch a feasibility study for a high-speed rail line in the western Taiwan corridor, which was completed in 1990. The study found that in a comparison of potential solutions to traffic problems in the corridor, a high-speed rail line would offer the highest transit volume, lowest land use, highest energy savings, and least pollution. In July 1990 the Preparation Office of High Speed Rail (POHSR) was established and a route was selected in 1991. Plans for the THSR were subsequently approved by the Executive Yuan in June 1992 and by Taiwan's legislature, the Legislative Yuan, in 1993.
Build-Operate-Transfer
In November 1994, Taiwan passed a law regarding the use of private finance in infrastructure projects, which also applied to the up-to-then state-run THSR project. Consequently, in 1995, POHSR was transformed into the Bureau of High Speed Rail (BOHSR), which started to tender THSR as a build-operate-transfer (BOT) scheme in October 1996.
The bidding process pitted Taiwan High Speed Rail Consortium (THSRC) against the Chunghwa High Speed Rail Consortium (CHSRC). THSRC's bid was based on the high-speed technology platform of Eurotrain, a joint venture between GEC-Alsthom, the main maker of the French TGV, and Siemens, the main maker of the German ICE, while CHSRC's bid was based on Japanese Shinkansen technology supplied by Taiwan Shinkansen Consortium (TSC), a joint venture of Japanese companies. THSRC, which submitted the lower bid and promised to build the line with zero net cost from the government, was chosen as the preferred bidder in September 1997. The group was renamed and formally established as the Taiwan High Speed Rail Corporation (THSRC) in May 1998. THSRC and the government signed the BOT agreement on 23 July 1998.
However, controversy arose during rolling-stock selection. In May 1999, as THSRC faced difficulties in raising capital, the government of Japan promised soft loans if THSRC switched to TSC. Although Eurotrain promised to match TSC's financial proposal, the Eschede train disaster in combination with TSC offering the newer 700 Series Shinkansen, convinced THSRC to reopen its core system bid, ultimately resulting in TSC selected as the preferred rolling-stock supplier in December 1999. Although Eurotrain eventually conceded in the bid, in February 2001 it filed for a US$800 million damage claim against THSRC at the Singapore International Arbitration Centre. After a lengthy arbitration process, the court ruled in March 2004 that THSRC should pay a compensation for the US$32.4 million Eurotrain spent on development and US$35.7 million for unjust enrichment. THSRC agreed to pay US$65 million (US$89 million with interest) to Eurotrain in November 2004.
Construction
Construction of the line by THSRC officially started in March 1999. Tunnels and other major civil engineering works were completed by 2004, along with the first delivery of the 700T trains. Testing and commissioning of the line then took place in 2005 and 2006, with a maximum testing speed of achieved in October 2005.
Opening
The railway was opened in 2007, with limited commercial services between Banqiao and Zuoying stations from 5 January, with full service from Taipei Station to Kaohsiung from May 2007.
Three additional stations located along the line – Miaoli, Changhua and Yunlin – opened in 2015.
Future plans
Southern extension
On 10 September 2019, the Executive Yuan announced that the railway would be expanded to Pingtung. Out of four proposed route options, it was confirmed on 27 September that the expansion would bypass central Kaohsiung, branching from Zuoying east towards western Pingtung City, near , with an estimated cost of NT$55.4 billion. Although lowest in cost, the option was met with criticism regarding its economic benefits.
The extension to Pingtung was approved by Premier Su Tseng-chang in January 2023, with opening of the extension planned for 2029.
On 28 December 2024, Executive Yuan announced that the extension route would be altered to pass through the city centre with a stop at Kaohsiung Main Station.
Northern extension
The line was extended from Taipei to Nangang, opening in July 2016.
On 25 October 2019, the Railway Bureau published an assessment report to extend the line further from Taipei to Yilan, cutting travel time to 13 minutes. The extension was approved in October 2020, and is planned to open by 2030.
Rolling stock
THSR 700T
Taiwan High Speed Rail started operation with 30 THSR 700T trainsets supplied by a consortium led by Kawasaki Heavy Industries. In response to increasing ridership and new stations that would begin operation in 2015, THSRC signed the contract for four new 700T trainsets with the Kawasaki consortium in May 2012 in Tokyo, Japan. The first (TR 31) trainset was delivered to Taiwan on 23 December 2012; the second (TR 32) on 21 January 2013; the third (TR33) on 25 January 2014; the fourth (TR34) on 12 August 2015.
The THSR 700T trainset is based on the 700 Series Shinkansen trainset used by JR Central and JR West in Japan. This marked the first time Shinkansen technology was exported to a foreign country, and it involved "rolling stock derived from a JR Central design running on both the European and Japanese track systems." Customization was focused on adapting to Taiwan's climate and geography, and the nose shape was optimized for tunnels wider than those in Japan.
The maximum service speed of the trains was raised from the 700 Series Shinkansen's . The 12 cars of a 700T train are grouped in three traction units with three power cars and one trailer each, providing of power; both end cars are trailers to avoid slip on powered bogies. The train is long and has a mass of when empty. The trains have a passenger capacity of 989 seats in two classes: 66 seats in 2+2 configuration in the single Business Car and 923 seats in 2+3 configuration in the eleven Standard Cars. The per capita energy consumption of a fully loaded 700T train is 16% of that of private cars and half that of buses; carbon dioxide emissions are 11% of private cars and a quarter that of buses.
N700S
In the late 2010s, THSRC began work to purchase additional high speed trains, in light of growing demand. 12 trains would be ordered, at a cost of around NT$30 billion. Due to the limited number of Japanese companies who build Shinkansen rolling stock, it took several years for THSRC to agree an acceptable deal. In 2022, it was reported that THSRC was speaking to European train manufacturers instead, as the price offered by Japanese companies was "unreasonable".
In March 2023, it was announced that a joint bid by Hitachi and Toshiba had been awarded the contract. Twelve of the latest generation of Shinkansen train—the N700S Series—would be delivered at a cost of NT$28 billion.
Engineering trains
THSRC uses overhead line inspection trains from Windhoff, Harsco railgrinders, Plasser & Theurer track tampers, and several ex-JR rolling stock to maintain its line. Among the latter include the JNR Class DD14 and JNR Class DD16 diesel-hydraulic locomotives, which were originally used for snowploughing by JR. The two ex-JR locomotives with THSRC are equipped with Shinkansen-style rotary couplers and standard-gauge bogies instead of the original gauge bogies and knuckle couplers and are used for shunting the 700T trainsets within the depot. THSRC also uses a former 0 Series Shinkansen end car as a structure gauge test car.
Operation
As the first high speed railway system in Taiwan, THSRC started operation in 2007 with many foreign employees, including French and German train drivers and operation controllers in the Operation Control Center (OCC). At the same time, THSRC also started to train local drivers and controllers. Since May 2008, all controllers working in the OCC have been Taiwanese, and since October 2008, all train drivers have been Taiwanese.
The OCC's main responsibility is to maintain safe train operations. THSRC has 132 controllers (July 2012), of which about one quarter are female, working 24 hours per day and 365 days per year in the OCC. Requirements for becoming a Chief Controller (主任控制員) include experience in all nine OCC positions, 300-hours of on-the-job training and acquiring qualification.
THSRC has 144 drivers (July 2012), of which almost 10% are female. All driver candidates must spend 8 months completing 1,326 hours of professional training and pass the National Certification before they can drive the train. In addition, after becoming a certified high-speed train driver, they undergo further on-the-job training at least three times each year in order to guarantee they can drive the train safely.
Natural disasters
Taiwan frequently faces multiple types of natural disasters, including typhoons, earthquakes, heavy rainfall, floods, and landslides. For this reason, a primary focus of THSRC's infrastructure design was how to respond to natural disasters such as earthquakes and how to ensure safety for all passengers and trains in any emergency situation.
THSRC has established a system to respond to natural disasters and unexpected intrusion onto the right-of-way, called DWS (Disaster Warning Systems). This system consists of a network of sensors installed along the rail route to detect unexpected situations such as earthquakes, strong winds, heavy rainfall, floods, landslides, and intrusions. In case of an unexpected situation, the DWS will send signals to the OCC (Operation Control Center) immediately; it will also activate contingency measures to ensure the safety of the passengers and trains, including decelerating or stopping trains in the affected area.
The DWS has functioned successfully since its initial operation in 2007. The most powerful earthquake that THSRC has experienced measured 6.4 on the Richter Scale with an epicenter from Jiaxian, Kaohsiung that shook southern Taiwan on 4 March 2010 (甲仙地震). One operating train was slightly derailed in Sinshih, Tainan, and six trains were stopped on the track. In spite of the temporary suspension of operations, there was no damage or casualties. All 2,500 affected passengers were evacuated in two hours without injury. Service resumed the next day. Such a record was well noted, and provided valuable experience in operational safety to the global railway industry.
In April 2010, it was reported that subsidence had been observed during construction on a viaduct section in Yunlin County. The subsidence continued, reaching up to over seven years. By 2010 subsidence had slowed, which was ascribed to the closure of some deep groundwater wells operating in the region. Although the situation was deemed safe with differential settlement between adjacent piers along the viaduct at only a sixth of the permissible level, the BOHSR urged the closure of more wells. On 25 July 2011, the government announced plans to close almost 1,000 wells in Changhua and Yunlin counties, reducing the amount of water pumped from deep wells by by 2021.
Service
According to THSR's July 2018 timetable, there are 989 train services per week of operation, with operation times between 05:50 to 24:00 every day. Most southbound trains originate from Nangang station and most northbound trains originate from Zuoying; however, a few trains operate just between Nangang and Taichung or between Taichung and Zuoying. Southbound trains are designated by odd train numbers, and northbound trains by even train numbers.
Each train consists of 1 business car (car 6) and 11 standard cars (including reserved seats and non-reserved seats). Since July 2010, non-reserved seats are available in cars 10 through 12 (some trains available in cars 9 through 12 or available in cars 8 through 12 ). Car 7 of each train is fitted with 4 wheelchair accessible chairs and a disabled-friendly restroom. Passengers can call THSR's Customer Service Hot Line at (Taiwan) 4066-3000 or visit any THSR station ticket window to reserve these seats.
By August 2012, implementation of 4G WiMAX on-board trains is expected to provide smooth wireless broadband services, making THSR the first high-speed ground transportation system equipped with this service.
In 2012, THSRC rated highly in the CommonWealth Magazine (天下雜誌) "Golden Service Award survey" (金牌服務大賞), not only far outpacing all rivals in the "long-distance land transport" category, but also taking the top spot in the overall rankings of 300 industries.
Local connections
To improve local public transit connections to THSR stations, the TRA built two new spur lines branching off from West Coast Line.
Shalun Line for Tainan opened on 2 January 2011,
Liujia Line for Hsinchu opened on 11 November 2011.
Stop patterns
With a few exceptions, the services follow the below pattern.
Ticket fare and discount
As of January 2018, a one-way Taipei–Zuoying trip, a THSR standard car adult ticket is NT$1490, and a business car ticket fare is NT$1950. The cost of a non-reserved seat is approximately 3% less than the regular price. Business and standard car reserved ticket reservations are available 28 days prior to the date of departure (including the departure day).
Senior citizens (Taiwan citizens above 65 years of age), registered disabled persons plus one accompanying passenger (Taiwan citizens only), and children (passengers under 12 years of age) are eligible for concession (half price) tickets.
A group discount is offered for groups of 11 or more. A group discount cannot be used in combination with other discount offers and does not include non-reserved seats. Passengers eligible for multiple discounts can only choose one discount offer.
Since 1 July 2010, a smart card system has provided frequent travelers with multi-ride (eight trips) or periodic tickets. THSR's contact-less smart cards allow the cardholder to travel between specific stations within a given time period for a certain number of rides. The card is sold in either registered (name-bearing) or non-registered form. Only adult tickets are available in this format, and cannot be used for rides between Banqiao and Taipei.
After purchasing or adding value to a multi-ride card, the card balance is valid for 45 days counted from the day of first use. The ticket is good for 8 rides. The multi-ride card provides a discount of about 21% off the full fare of a reserved Standard Seat. Non-registered and registered multi-ride tickets can be purchased at the ticket windows of all THSRC stations. Upon first purchase of a multi-ride ticket, a card deposit fee of NT$100 is required (refundable if the card is returned). The registered multi-ride ticket is limited to personal use by the registered cardholder.
Since November 2012, an Early Bird discount of 35% has been offered for a limited number of tickets sold no later than 8 days before the departure date. If the 35% off tickets sell out before the deadline, tickets with a discount of 20% off are offered. If these tickets sell out before the deadline, tickets with a discount of 10% off are offered. If all early bird tickets are sold out, then full fare tickets are offered.
Train frequency
THSRC operates additional train services during national holidays. On 29 June 2011, a proposal by THSRC to increase the maximum number of train services to 210 per day (compared to the existing 175 per day) passed an environmental impact assessment, increasing the number of possible services on "high-load days".
Ridership
Original estimates predicted a daily ridership of 180,000 after launch, growing to 400,000 by 2036. In view of a 50% drop in airline passengers in the wake of the 1997 Asian financial crisis, forecasts were revised downwards. The final initial ridership estimate was 140,000 passengers per day. Actual initial ridership did not match these projections. In September 2007, six months after opening, THSRC carried 1.5 million passengers monthly, translating to about 50,000 passengers daily. In the second year, passenger numbers almost doubled. In the third year, average daily ridership continued to grow to 88,000 passengers per day, jumping to over 120,000 passengers per day in 2012. (updated to September 2012) Seat occupancy was around 45% in the first three years, with a modest improvement achieved in 2009, and reached 53.91% in 2012. (updated to September 2012) Punctuality is stable above 99%.
The 10-millionth passenger was carried after 265 days of operation on 26 September 2007, while the 100-millionth passenger was carried after 1,307 days on 3 August 2010, and 200-millionth by December 2012. On 10 October 2011, the Double Ten Day holiday, THSRC transported a single-day record of 189,386 passengers. On 5 February 2011, the third day of Chinese New Year’s celebration, a new record of 190,596 passengers was achieved. The next single-day record was reached on 25 January 2012, also the third day of Chinese New Year's celebration, at 191,989 passengers. The most recent record is 212,000 passengers transported on 1 January 2013.
The high-speed trains have successfully out-competed planes: by August 2008, half of the air routes between Taipei and the country's western cities had been discontinued, including all connections between cities with THSR stations except for a single daily connection between Taipei and Kaohsiung. Total domestic air traffic was expected to be halved from 2006 to 2008, and actually fell from 8.6 to 4.9 million. In June 2012, officials announced the discontinuation of the last remaining commercial flight between Taipei and Kaohsiung. The share for conventional rail between Taipei and Kaohsiung fell from 9.71% in 2006 to 2.5% in 2008, while high-speed rail became the most common mode of transport at 50% of all trips by 2008. The opening of THSR led to a 10% reduction of traffic on the parallel expressway in 2007. Despite cheaper ticket prices, long-distance bus companies reported that passenger volumes had fallen by 20 to 30 percent by 2008.
Infrastructure
Construction of the system took more than 2,000 professional engineers from 20 countries and over 20,000 foreign and domestic workers six years to complete. Construction work was broken into several specialized lots that were contracted separately. One group of contracts was for civil works, covering the construction of the superstructure of open line sections. Stations and depots were the subject of separate groups of construction contracts. A fourth group of contracts was for track work.
The Taiwan North-South High Speed Rail Project was awarded the first prize for the Outstanding Civil Engineering Project Award by the Asian Civil Engineering Coordination Council (ACECC) in Sydney in 2010.
In 2011, the Public Construction Commission (公共工程委員會) organized an on-line voting campaign that garnered over 330,000 votes, to select the 100 best infrastructure projects (百大建設) in Taiwan to celebrate the centennial of the Republic; Taiwan High Speed Rail topped the list.
Track
Reflecting a design speed of , track layout was designed with a minimum curve radius of , track-centre distance of , right-of-way width of , and a maximum gradient of 2.5%, except for 3.5% at one location. All but of track is ballastless, combining slab track of Japanese manufacture on open line sections with switches from a German supplier. Track laying began in July 2003. The line was electrified with the 25 kV/60 Hz AC system. The signalling and train control system was laid out for bi-directional operation according to European specifications. Each track section has a checkpoint, and an automatic control system ensures that trains are spaced at least apart to prevent collisions.
Most of the line is elevated. About or 73% of the line runs on viaducts, mostly precast pre-stressed concrete box girder spans, the first of which was put in place in October 2001.
The Changhua-Kaohsiung Viaduct is a continuous section from Baguashan (八卦山) in Changhua County to Zuoying in Kaohsiung. It was the second longest bridge in the world as of 2017. Viaducts were designed to be earthquake resistant to allow for trains to stop safely during a seismic event and for repairable damage following a maximum design earthquake. Bridges built over known fault lines were designed to survive fault movements without catastrophic damage.
About or 18% of the line is in tunnels, including of the TRUPO section in Taipei, as well as 48 tunnels with a total length of on the other sections, the longest of which is Paghuashan Tunnel, at a finished length of . Forty-two of the tunnels included a total of of mined sections, all of which were bored with the sequential excavation and support construction method, with excavated tunnel faces of , between November 2000 and July 2003. The finished interior cross-sectional area of , set according to wider European standards, provides space for two tracks with safety walkways.
After four months of delays, trial runs using the first THSR 700T trains began on 27 January 2005, on the Tainan–Kaohsiung section. On 30 October 2005, a day after a test run passed the planned top service speed of , the targeted maximum test speed of was achieved. The section between Banqiao (Taipei) and Zuoying (Kaohsiung) opened to the public on 5 January 2007. The HSR platforms at Taipei Station opened on 2 March 2007, bringing the entire line into operation.
Stations
A distinctive feature of the system's station placement is that many are located at the periphery of urban areas, rather than within city centers. The decision was made with the expectation that the stations would act as centers for planned communities and thus increase the property values of the surrounding area. A study in 2010 showed that this isn't the case, but later analyses show that property prices around certain stations have indeed risen. Since the THSR's opening, cities have gradually expanded their mass transit systems to connect with these stations.
Environmental issues
Environmental mitigation measures in the line's construction phase included the construction of animal bridges over the line, the planting and re-planting of trees along the track as noise screens, and the purchase of farmland to create a preservation area for jacana birds away from the line.
THSRC is involved in the preservation of the pheasant-tailed jacana, which is considered endangered in Taiwan. An artificial habitat recovery project was completed in collaboration with the local government, country development organizations and non-profit organizations for a cost of NT$50 million. In 2007, the recovery habitat was officially renamed the "Pheasant-tailed Jacana Eco-Educational Nature Park" and since then, it has opened to the public. THSRC arranges for elementary and junior high school students to visit the park annually.
A 330 year old camphor tree and a temple in Hsinchu County are located on the main route of the THSR, and both of them faced removal because of railway construction. The temple established beside the old tree serves as a major religious site for the local community. In 1998, THSRC adjusted the line and design to keep the tree and temple in their original place and cooperated with the local government and people to protect the old tree and the temple until today. Afterwards, together with the local government, the Environment and the Resources Protection Committee, and cultural and historical authorities, THSRC drafted the Hsinchu Old Camphor Tree Medical Plan, which called for the repair of decayed branches as well as measures designed to maintain the long-term growth and the health of the tree.
Financial
Revenue and cost
Most of THSRC's revenue comes from ticket sales; supplemental income comes from other activities such as advertising and renting spaces for standing shops and spots in plazas. Advertising spots on trains and station platforms have also been sold.
Revenues grew along with ridership over the first three years, but ridership remained below expectations. In 2008 the second year of operation, revenues fell barely short of THSRC's expectations a year earlier of a doubling of first-year results.
The cost of running the trains and infrastructure, or cash operating costs, was initially over NT$1 billion a month, but was reduced to around NT$850-900 million in 2008. Revenues first exceeded this level, thus generating a positive operating cash flow, in the fourth month of operation (April 2007).
For THSRC, the over heavy accounting of the fixed cost of fixed assets like rolling stock and infrastructure (depreciation) is a significant non-cash element of total operating costs. In its first two years of operation, THSRC applied straight-line depreciation, distributing costs evenly over a period of 26.5 years. As a result, the balance of operating revenues and costs (operating income) showed a high loss in the first year of operation, which was only reduced as revenues grew in the second year. The depreciation period set for THSRC reflected the length of the B.O.T. concession rather than the much longer lifespan of the infrastructure, and it is the factor for the operating loss. After adopting an activity depreciation method which is variable in time, THSRC posted its first operating profit for 2009, the third year of operation. The company reported its first annual profit of NT$5.78 billion after five years of operation.
For the first time in its five-year operation, the Company reported a net income of NT$5.78 billion, with earnings per share of NT$0.59. Revenues increased by 16.65% from NT$27.64 billion to NT$32.24 billion, with operating costs and expenses (excluding depreciation and amortization) rising by only 4.98%. Over the same period of time, gross profit totaled NT$12.98 billion (an increase of 30.32%), income from operations totaled NT$12.06 billion (an increase of 32.93%) and EBITDA totaled NT$22.73 billion (an increase of 22.34%). 2011 gross profit, income from operation and EBITDA were all record highs. Since commencing operations in 2007, THSRC has had a significant influence on Taiwan's economy and on the lives of its people. In 2011, the Company continued to pursue sustainable growth aligned with the interests of shareholders and society, achieving record profits even amid a challenging economic environment.
The interest cost is another major item of this company's financing. In the first few years of operation, interest rates were well above market rates. Interest expense per month stood at around NT$1.3 billion in 2008, when THSRC first achieved break-even cash flow, with revenue and cash expenses (which exclude depreciation) both around NT$2.1 billion in 2008. Interest rates fell in the first half of 2009, reducing interest expenses and contributing to a reduced net loss.
In 2010, THSRC obtained a new syndicate loan to alleviate its imminent financial burden. The company signed a NT$382 billion refinancing contract with a consortium of eight domestic banks led by the Bank of Taiwan, and used the new loan to pay off the previous syndicated loan, which had higher interest. , the long-term debts totalling NT$385 billion included NT$26 billion in corporate bonds and NT$359 billion in bank loans. In comparison with the terms and conditions of previous loans, the refinancing debts carried lower interest rates and longer tenors, up to 22 years.
Financial and loan
In cumulative figures, until July 2008, depreciation and interest were equal to 95% of THSRC's accumulated debt. Both THSRC and a September 2009 government report identified an unreasonable financial structure and the resulting high interest rates and high depreciation charges as the main causes of negative financial performance, while the government assessed THSRC to have performed well in its core business, as measured by earnings before interest, taxes, depreciation and amortization (EBITDA). To reduce its interest load, THSRC sought to revise its loan structure in 2008 and again in 2009. To reduce depreciation costs by increasing the amortization time, THSRC requested an extension of its 35-year concession period.
By the summer of 2009, THSRC's cumulative losses were equivalent to two-thirds of its equity capital. In response to global financial crisis and domestic economic recession, THSRC proposed to increase income and reduce expenditures in several aspects in the hope of raising operation performances. In February 2009, THSRC announced to adjust train frequency, cut down salary payment by 10~20% among management level, and measured to expand fare promotion to stimulate ridership.
While the media questioned whether the planned construction of three more intermediate stations and the extension to Nangang would be postponed, THSRC published press release on 28 September 2009, stating that the company will comply with "Taiwan High Speed. Rail Construction and Operation Contract", and the construction project of 3 intermediate stations, namely Miaoli, Changhua and Yunlin will be initiated in July 2012, and is scheduled to start its operation from 2015. By the time of completion, there will be a total of 12 stations along the THSRC operation route. The company was put under new management in September 2009 with the aim of turning around the company's finances with government help in arranging refinancing of the loans.
The government took majority control of the company after the election of its new board on 10 November 2009. In January 2010, when accumulated losses already exceeded NT$70 billion, THSRC signed a government-guaranteed refinancing deal in which eight government-dominated banks provided NT$382 billion at lower interest rates and longer maturity. The government also approved the company's new variable depreciation charge.
Incidents
On 12 April 2013, suspicious luggage items were found inside the North bound train No. 616 toilet when it was heading towards THSR Hsinchu Station. The train was stopped at THSR Taoyuan Station and all of the passengers were evacuated. Later, it was determined the luggage contained an unidentified liquid in cans, alarm clock and white particulate matter. The items were dismantled by the bomb squad and taken for further investigation. Two KMT legislators, Hsu Hsin-ying and Lu Shiow-yen, were on board.
Part of the tracks near Tainan were badly damaged during the earthquake on 6 February 2016. All high-speed rail services south of Chiayi Station were suspended until 7 February 2016.
On 10 May 2017, a non-passenger carrying train traveled in the opposite direction of the track from Zuoying to Tainan for due to human negligence.
Due to the COVID-19 pandemic, the THSR, along with Taiwan Railway Administration and bus services nationwide, began to require all passengers to wear surgical masks as of 1 April. In addition, infrared sensors were set up at twelve stations to detect fevers, eating and drinking were prohibited on board the trains, trains and stations were disinfected more frequently, and the THSR cancelled all non-reserved seating tickets, which allowed for crowds of passengers to stand if no seats were available. It was reported that the switch to reserved seats only aimed to reduce crowding.
Public relations activities
THSRC conducts community engagement activities to raise its profile.
Since 2009, the company has organized an annual "Ride THSR and Join the Book Exhibition for Free" event to promote a national reading culture; school-age passengers from remote villages are given free admission to the Taipei International Book Exhibition and go there on a themed high-speed "reading train", which features a celebrity reading a book over the train's public address system.
Since 2010, along with World Vision Taiwan, THSRC has run a tuition fee assistance program for thousands of underprivileged children, to which passengers contribute.
Other events have been a cappella singers at stations; gift-giving to couples taking wedding photos at major stations; station tours for the public and experience-sharing with its fellow railway transportation operators; and in collaboration with non-profit organizations, thousands of free rides to underprivileged groups and families.
Students at primary, secondary and tertiary level learn about high-speed rail and THSRC at "THSR Camps", held in partnership with the Railway Cultural Society of Taiwan, the National Chiao Tung University Railway Research Society, and the China Youth Corps.
In popular culture
The first film to feature THSR prominently was the 2007 Taiwanese movie Summer's Tail, directed by Cheng Wen-tang ().
Railfan: Taiwan High Speed Rail, a 2007 train simulator video game developed jointly by Taiwanese company Actainment and Japanese company Ongakukan on the basis of the latter's Train Simulator series, featured real video and was the first Taiwanese game for Sony Computer Entertainment's PlayStation 3 system.
The "National Geographic" website chose travel by Taiwan's high speed train as the "Best winter trip” in 2013.
| Technology | High-speed rail | null |
229017 | https://en.wikipedia.org/wiki/Latissimus%20dorsi%20muscle | Latissimus dorsi muscle | The latissimus dorsi () is a large, flat muscle on the back that stretches to the sides, behind the arm, and is partly covered by the trapezius on the back near the midline.
The word latissimus dorsi (plural: latissimi dorsi) comes from Latin and means "broadest [muscle] of the back", from "latissimus" () and "dorsum" (). The pair of muscles are commonly known as "lats", especially among bodybuilders.
The latissimus dorsi is responsible for extension, adduction, transverse extension also known as horizontal abduction (or horizontal extension), flexion from an extended position, and (medial) internal rotation of the shoulder joint. It also has a synergistic role in extension and lateral flexion of the lumbar spine.
Due to bypassing the scapulothoracic joints and attaching directly to the spine, the actions the latissimi dorsi have on moving the arms can also influence the movement of the scapulae, such as their downward rotation during a pull up.
Structure
Variations
The number of dorsal vertebrae, to which it is attached, varies from four to eight; the number of costal attachments varies; muscle fibers may or may not reach the crest of the ilium.
A muscle slip, the axillary arch, varying from 7 to 10 cm in length, and from 5 to 15 mm in breadth, occasionally springs from the upper edge of the latissimus dorsi about the middle of the posterior fold of the axilla, and crosses the axilla in front of the axillary vessels and nerves, to join the under surface of the tendon of the pectoralis major, the coracobrachialis, or the fascia over the biceps brachii. This axillary arch crosses the axillary artery, just above the spot usually selected for the application of a ligature, and may mislead a surgeon. It is present in about 7% of the population and may be easily recognized by the transverse direction of its fibers. Guy et al. extensively described this muscular variant using MRI data and positively correlated its presence with symptoms of neurological impingement.
A fibrous slip usually passes from the upper border of the tendon of the Latissimus dorsi, near its insertion, to the long head of the triceps brachii. This is occasionally muscular, and is the representative of the dorsoepitrochlearis brachii of apes. This muscular form is found in ~5% of humans and is sometimes termed the latissimocondyloideus.
The latissimus dorsi crosses the inferior angle of the scapula. A study found that, of 100 cadavers dissected:
43% had "a substantial amount" of fibers in the latissimus dorsi originating from the scapula.
36% had few or no muscular fibers, but a "soft fibrous link" between the scapula and the latissimus dorsi
21% had little or no connecting tissue between the two structures.
Triangles
The lateral margin of the latissimus dorsi is separated below from the obliquus externus abdominis by a small triangular interval, the lumbar triangle of Petit, the base of which is formed by the iliac crest, and its floor by the obliquus internus abdominis.
Another triangle is situated behind the scapula. It is bounded above by the trapezius, below by the latissimus dorsi, and laterally by the vertebral border of the scapula; the floor is partly formed by the rhomboideus major. If the scapula is drawn forward by folding the arms across the chest, and the trunk bent forward, parts of the sixth and seventh ribs and the interspace between them become subcutaneous and available for auscultation. The space is therefore known as the triangle of auscultation.
The latissimus dorsi can be remembered best for insertion as "A Miss Between Two Majors". As the latissimus dorsi inserts into the floor of the intertubercular groove of the humerus it is surrounded by two major muscles. The teres major inserts medially on the medial lip of the intertubercular groove and the pectoralis major inserts laterally onto the lateral lip.
Nerve supply
The latissimus dorsi is innervated by the sixth, seventh, and eighth cervical nerves through the thoracodorsal (long subscapular) nerve. Electromyography suggests that it consists of six groups of muscle fibres that can be independently coordinated by the central nervous system.
Function
The latissimus dorsi assists in depression of the arm with the teres major and pectoralis major. It adducts, extends, and internally rotates the shoulder. When the arms are in a fixed overhead position, the latissimus dorsi pulls the trunk upward and forward.
It has a synergistic role in extension (posterior fibers) and lateral flexion (anterior fibers) of the lumbar spine, and assists as a muscle of both forced expiration (anterior fibers) and an accessory muscle of inspiration (posterior fibers).
Most latissimus dorsi exercises concurrently recruit the teres major, posterior fibres of the deltoid, long head of the triceps brachii, among numerous other stabilizing muscles. Compound exercises for the 'lats' typically involve elbow flexion and tend to recruit the biceps brachii, brachialis, and brachioradialis for this function. Depending on the line of pull, the trapezius muscles can be recruited as well; horizontal pulling motions such as rows recruit both latissimus dorsi and trapezius heavily.
Training
The power/size/strength of this muscle can be trained with a variety of different exercises. Some of these include:
Vertical pulling movements such as pull-downs and pull-ups (including chin-ups)
Horizontal pulling movements such as bent-over row, T-bar row and other rowing exercises
Shoulder extension movements with straight arms such as straight-arm lat pulldowns and pull-overs
Deadlift
Clinical significance
Tight latissimus dorsi has been shown to be a contributor to chronic shoulder pain and chronic back pain. Because the latissimus dorsi connects the spine to the humerus, tightness in this muscle can manifest as either sub-optimal glenohumeral joint (shoulder) function which leads to chronic pain or tendinitis in the tendinous fasciae connecting the latissimus dorsi to the thoracic and lumbar spine.
The latissimus dorsi is a potential source of muscle for breast reconstruction surgery after mastectomy (e.g., Mannu flap) or to correct pectoral hypoplastic defects such as Poland's syndrome. An absent or hypoplastic latissimus dorsi can be one of the associated symptoms of Poland's syndrome.
Cardiac support
For heart patients with low cardiac output and who are not candidates for cardiac transplantation, a procedure called cardiomyoplasty may support the failing heart. This procedure involves wrapping the latissimus dorsi muscles around the heart and electrostimulating them in synchrony with ventricular systole.
Injury
Injuries to the latissimus dorsi are rare. They occur disproportionately in baseball pitchers. Diagnosis can be achieved by visualization of the muscle and movement testing. MRI of the shoulder girdle will confirm the diagnosis. Muscle belly injuries are treated with rehabilitation while tendon avulsion injuries can be treated surgically, or with rehab. Regardless of treatment, patients tend to return to play without any functional losses.
Additional images
| Biology and health sciences | Human anatomy | Health |
229045 | https://en.wikipedia.org/wiki/Pager | Pager | A pager, also known as a beeper or bleeper, is a wireless telecommunications device that receives and displays alphanumeric or voice messages. One-way pagers can only receive messages, while response pagers and two-way pagers can also acknowledge, reply to, and originate messages using an internal transmitter.
Pagers operate as part of a paging system which includes one or more fixed transmitters (or in the case of response pagers and two-way pagers, one or more base stations), as well as a number of pagers carried by mobile users. These systems can range from a restaurant system with a single low power transmitter, to a nationwide system with thousands of high-power base stations.
Pagers were developed in the 1950s and 1960s, and became widely used by the 1980s through the late 1990s and early 2000s. Later in the 21st century, the widespread availability of cellphones and smartphones with text messaging capability has greatly diminished the pager industry. Nevertheless, pagers continue to be used by some emergency services and public safety personnel, because modern pager systems' coverage overlap, combined with use of satellite communications, can make paging systems more reliable than terrestrial based cellular networks in some cases, including during natural and human-made disasters. This resilience has led public safety agencies to adopt pagers over cellular and other commercial services for critical messaging.
History
The first telephone pager system was patented in 1949 by Al Gross. Intended for the use of physicians, there was initial resistance to the idea of being permanently on-call, according to Gross.
One of the first practical paging services was launched in 1950 for physicians in the New York City area. Physicians paid US$12 per month and carried a pager that would receive phone messages within of a single transmitter tower. The system was manufactured by the Reevesound Company and operated by Telanswerphone. In 1960, John Francis Mitchell combined elements of Motorola's walkie-talkie and automobile radio technologies to create the first transistorized pager, and from that time, paging technology continued to advance and pager adoption among emergency personnel was still popular as of July 2016.
In 1962, the Bell System, the U.S. telephone monopoly, presented its Bellboy radio paging system at the Seattle World's Fair. Bellboy was the first commercial system for personal paging. It also marked one of the first consumer applications of the transistor (invented by Bell Labs in 1947), for which three Bell Labs inventors received a Nobel Prize in Physics in 1956. Solid-state circuitry enabled the Bellboy pager, about the size of a small TV remote device, to fit into a customer's pocket or purse, quite a feat at that time. The Bellboy was a terminal that notified the user when someone was trying to call them. Bell System Bellboy radio pagers each used three reed receiver relays, each relay tuned to one of 33 different frequencies, selectively ringing a particular customer when all three relays were activated at the same time—a precursor of DTMF. When the person received an audible signal (a buzz) on the pager, the user found a telephone and called the service center, which informed the user of the caller's message.
In the mid-1980s, tone and voice radio paging became popular among emergency responders and professionals. Tone and voice pagers were activated either by a local base station, or through a telephone number assigned to each individual pager. In the 1990s, pagers became popular among the general public as a cheaper, smaller, and more reliable alternative to mobile phones. The ReFLEX protocol was developed in the mid-1990s.
As prices for mobile phones declined, small form factor phones like the Motorola StarTAC and the Nokia Series 40 line came on the market, cellular connectivity expanded, and digital phones adopted text messaging, most pager customers outside of specialist fields migrated to mobile phones toward the end of the 1990s. While Motorola announced the end of its new pager manufacturing in 2001, pagers remained in use in large hospital complexes. First responders in rural areas with inadequate cellular coverage are often issued pagers.
The 2005 London bombings resulted in overload of TETRA systems by the emergency services and showed that pagers, with their absence of necessity to transmit an acknowledgement before showing the message, and the related capability to operate on very low signal levels, are not completely outclassed by their successors. Volunteer firefighters, EMS paramedics and rescue squad members usually carry pagers to alert them of emergency call outs for their department. These pagers receive a special tone from a fire department radio frequency.
Restaurant pagers remain in wide use since the 2000s. Customers are given a portable receiver that would usually vibrate, flash, or beep when a table becomes free or when their meal is ready. Pagers have been popular with birdwatchers in Great Britain and Ireland since 1991, with companies Rare Bird Alert and Birdnet Information offering news of rare birds sent to pagers that they sell.
Today, companies like Visiplex offer similar solutions for onsite pager systems in the medical, education and commercial sectors.
Decline
By early 2002, pager usage was rapidly declining in places like North America due to the proliferation of cellular telephones.
The U.S. paging industry generated $2.1 billion in revenue in 2008, down from $6.2 billion in 2003. In Canada, 161,500 Canadians paid $18.5 million for pager service in 2013. Telus, one of the three major mobile carriers, announced the end to its Canadian pager service as of 31 March 2015, but rivals Bell, Rogers and PageNet intend to continue service.
In 2017 the UK National Health Service was thought to have been using over 10% of the remaining pagers in the world (130,000), with an annual cost of £6.6 million. Matt Hancock, (then) Secretary of State for Health and Social Care, announced in February 2019 that the 130,000 pagers still in use were to be phased out. NHSX announced plans in May 2020 to replace pagers and bleepers with "more modern communication tools," accelerated by the pressure placed on the service by the COVID-19 pandemic in England. In August 2020, a new procurement framework for clinical communications was launched which was intended to phase out pagers by the end of 2021, replacing them with "dedicated clinical-facing communication and tasks management tools" from 25 approved suppliers.
In Japan, more than ten million pagers were active in 1996. On 1 October 2019, Japan's last paging service provider shut down radio signals and terminated its service.
In Russia, the last paging provider was closed in November 2021.
Design
Many paging network operators now allow numeric and textual pages to be submitted to the paging networks via email. A significant convenience for users given the widespread adoption of email, and commonalities in delivery assurances. This can result in pager messages being delayed or lost. Older forms of message submission using the Telelocator Alphanumeric Protocol involve modem connections directly to a paging network and are less subject to these delays. For this reason, older forms of message submission retain their usefulness for disseminating highly-important alerts to users such as emergency services personnel.
Common paging protocols include TAP, FLEX, ReFLEX, POCSAG, GOLAY, ERMES and NTT. Past paging protocols include Two-tone and 5/6-tone. In the United States, pagers typically receive signals using the FLEX protocol in the 900 MHz band. Commercial paging transmitters typically radiate 1000 watts of effective power, resulting in a much wider coverage area per tower than a mobile phone transmitter, which typically radiates around 0.6 watts per channel. Although FLEX paging networks tend to have stronger in-building coverage than mobile phone networks, commercial paging service providers will work with large institutions to install repeater equipment in the event that service is not available in needed areas of the subscribing institution's buildings. This is especially critical in hospital settings where emergency staff must be able to reliably receive pages to respond to patient needs.
Unlike mobile phones, most one-way pagers do not display any information about whether a signal is being received or about the strength of the received signal. Since one-way pagers do not contain transmitters, one-way paging networks have no way to track whether a message has been successfully delivered to a pager. Because of this, if a one-way pager is turned off or is not receiving a usable signal at the time a message is transmitted, the message will not be received and the sender of the message will not be notified of this fact. In the mid-1990s, some paging companies began offering a service, which allowed a customer to call their pager number and have numeric messages read back to them. This was useful for times when the pager was off or out of the coverage area, as it would know what pages were sent to the subscriber even if the subscriber never actually received the page. Other radio bands used for pagers include the 400 MHz band, the VHF band and the FM commercial broadcast band (88–108&MHz). Other paging protocols used in the VHF, 400 MHz UHF and 900 MHz bands include POCSAG and ERMES. In Canada and the United States, pagers that use the commercial FM band receive a subcarrier, called the Subsidiary Communications Authority, of a broadcast station. On-site paging systems in hospitals, unlike wide area paging systems, are local area services. Hospitals commonly use on-site paging for communication with staff and increasingly for contacting waiting patients when their appointment is due. These offer waiting patients the opportunity to leave the waiting area, but still be contacted.
Operation
Paging systems are operated by commercial carriers, often as a subscription service and they are also operated directly by end users as private systems. Commercial carrier systems tend to cover a larger geographical area than private systems, while private systems tend to cover their limited area more thoroughly and deliver messages faster than commercial systems. In all systems, clients send messages to pagers, an activity commonly referred to as paging. System operators often assign unique phone numbers or email addresses to pagers (and pre-defined groups of pagers), enabling clients to page by telephone call, e-mail and SMS. Paging systems also support various types of direct connection protocols, which sacrifice global addressing and accessibility for a dedicated communications link. Automated monitoring and escalation software clients, often used in hospitals, IT departments and alarm companies, tend to prefer direct connections because of the increased reliability. Small paging systems, such as those used in restaurant and retail establishments, often integrate a keyboard and paging system into a single box, reducing both cost and complexity.
Paging systems support several popular direct connection protocols, including TAP, TNPP, SNPP and WCTP, as well as proprietary modem- and socket-based protocols. Additionally, organizations often integrate paging systems with their Voice-mail and PBX systems, conceptually attaching pagers to a telephone extension and set up web portals to integrate pagers into other parts of their enterprise. A paging system alerts a pager (or group of pagers) by transmitting information over an RF channel, including an address and message information. This information is formatted using a paging protocol, such as 2-tone, 5/6-tone, GOLAY, POCSAG, FLEX, ERMES, or NTT. Two-way pagers and response pagers typically use the ReFLEX protocol.
Modern paging systems typically use multiple base transmitters to modulate the same signal on the same RF channel, a design approach called simulcast. This type of design enables pagers to select the strongest signal from several candidate transmitters using FM capture, thereby improving overall system performance. Simulcast systems often use satellite to distribute identical information to multiple transmitters and GPS at each transmitter to precisely time its modulation relative to other transmitters. The coverage overlap, combined with use of satellite communications, can make paging systems more reliable than terrestrial based cellular networks in some cases, including during natural and human-made disaster. This resilience has led public safety agencies to adopt pagers over cellular and other commercial services for critical messaging.
Categories
Pagers themselves vary from very cheap and simple beepers, to more complex personal communications equipment, falling into eight main categories.
Beepers or tone-only pagers Beepers or tone-only pagers are the simplest and least expensive form of paging. They were named beepers because they originally made a beeping noise, but current pagers in this category use other forms of alert as well. Some use audio signals, others light up and some vibrate, often used in combination. The majority of restaurant pagers fall into this category.
Voice/tone Voice/Tone pagers enable pager users to listen to a recorded voice message when an alert is received.
Numeric Numeric Pagers contain a numeric LCD display capable of displaying the calling phone number or other numeric information generally up to 10 digits. The display can also convey pager codes, a set of number codes corresponding to mutually understood pre-defined messages.
Alphanumeric Alphanumeric pagers contain a more sophisticated LCD capable of displaying text and icons. These devices receive text messages, often through email or direct connection to the paging system. The sender must enter a message, either numeric and push # or, text & push # or a verbal message. The pager does not automatically record the sender's number; the pager will beep but no message can be seen or heard if none has been entered.
Response Response pagers are alphanumeric pagers equipped with built-in transmitters, with the ability to acknowledge/confirm messages. They also allow the user to reply to messages by way of a multiple-choice response list, and to initiate "canned" messages from pre-programmed address and message lists. These devices are sometimes called "1.5-way pagers" or "1.7-way pagers" depending on capabilities.
Two-way Two-way pagers are response pagers with built-in QWERTY keyboards. These pagers allow the user to reply to messages, originate messages and forward messages using free-form text as well as "canned" responses.
One-way modems One-way modems are controllers with integrated paging receivers, which are capable of taking local action based on messages and data they receive.
Two-way modems Two-way modems have capabilities similar to one-way modems. They can also confirm messages and transmit their own messages and data.
Security
Pagers have certain privacy advantages and disadvantages compared with cellular phones. Since a one-way pager is a passive receiver only (it sends no information back to the base station), its location cannot be tracked. However, this can also be disadvantageous, as a message sent to a pager must be broadcast from every paging transmitter in the pager's service area. Thus, if a pager has nationwide service, a message sent to it could be intercepted by criminals or law enforcement agencies anywhere within the nationwide service area.
Attacks
On September 17, 2024, a massive attack against Hezbollah members in Lebanon and Syria was allegedly committed by Israel, who simultaneously detonated pagers that they were using. Lebanese Health authorities confirmed at least nine deaths and over 3,000 injuries as a result of the explosions. In February 2024, Hezbollah leader Hassan Nasrallah had told the group's members to use pagers instead of cell phones, claiming that Israel had infiltrated their mobile phone network.
In popular culture
As is the case with many new technologies, the functionality of the pager shifted from necessary professional use to a social tool integrated in one's personal life. During the rise of the pager, it became the subject of various forms of media, most notably in the 1990s hip-hop scene. Popular artists from the era, including Ice Cube, Method Man, and A Tribe Called Quest, began referencing newly developed mobile technologies such as the pager. A Tribe Called Quest's single "Skypager" directly speaks of the importance of such a wireless communication device, with group member Q-Tip stating that the Skypager "serves an important communicative function for a young professional with a full calendar". Three 6 Mafia's "2-Way Freak," Sir Mix-A-Lot's "Beepers" and "Bug a Boo" from Destiny's Child also make reference to pagers.
Illicit drug dealers used pagers to great effect during the 1990s to conduct commerce, using them to arrange meetings with buyers. Associate superintendent for Miami-Dade County Public Schools in Florida Gul James Fleming once called them "the most dominant symbol of the drug trade" and schools have previously forbidden students from carrying them because of the ease with which they could be "used to arrange illegal drug sales."
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229063 | https://en.wikipedia.org/wiki/Schizoid%20personality%20disorder | Schizoid personality disorder | Schizoid personality disorder (, often abbreviated as SzPD or ScPD) is a personality disorder characterized by a lack of interest in social relationships, a tendency toward a solitary or sheltered lifestyle, secretiveness, emotional coldness, detachment, and apathy. Affected individuals may be unable to form intimate attachments to others and simultaneously possess a rich and elaborate but exclusively internal fantasy world. Other associated features include stilted speech, a lack of deriving enjoyment from most activities, feeling as though one is an "observer" rather than a participant in life, an inability to tolerate emotional expectations of others, apparent indifference when praised or criticized, all forms of asexuality, and idiosyncratic moral or political beliefs.
Symptoms typically start in late childhood or adolescence. The cause of SzPD is uncertain, but there is some evidence of links and shared genetic risk between SzPD, other cluster A personality disorders, and schizophrenia. Thus, SzPD is considered to be a "schizophrenia-like personality disorder". It is diagnosed by clinical observation, and it can be very difficult to distinguish SzPD from other mental disorders or conditions (such as autism spectrum disorder, with which it may sometimes overlap).
The effectiveness of psychotherapeutic and pharmacological treatments for the disorder has yet to be empirically and systematically investigated. This is largely because people with SzPD rarely seek treatment for their condition. Originally, low doses of atypical antipsychotics were used to treat some symptoms of SzPD, but their use is no longer recommended. The substituted amphetamine bupropion may be used to treat associated anhedonia. However, it is not general practice to treat SzPD with medications, other than for the short-term treatment of acute co-occurring disorders (e.g. depression). Talk therapies such as cognitive behavioral therapy (CBT) may not be effective, because people with SzPD may have a hard time forming a good working relationship with a therapist.
SzPD is a poorly studied disorder, and there is little clinical data on SzPD because it is rarely encountered in clinical settings. Studies have generally reported a prevalence of less than 1%. It is more commonly diagnosed in males than in females. SzPD is linked to negative outcomes, including a significantly compromised quality of life, reduced overall functioning even after 15 years, and one of the lowest levels of "life success" of all personality disorders (measured as "status, wealth and successful relationships"). Bullying is particularly common towards schizoid individuals. Suicide may be a running mental theme for schizoid individuals, though they are not likely to attempt it. Some symptoms of SzPD (e.g. solitary lifestyle, emotional detachment, loneliness, and impaired communication), however, have been stated as general risk factors for serious suicidal behavior.
History
The term schizoid was coined in 1908 by Eugen Bleuler to describe a human tendency to direct attention toward one's inner life and away from the external world. Bleuler labeled the exaggeration of this tendency the "schizoid personality". He described these personalities as "comfortably dull and at the same time sensitive, people who in a narrow manner pursue vague purposes". In 1910, August Hoch introduced a very similar concept called the "shut-in" personality. Characteristics of it were reticence, reclusiveness, shyness and a preference for living in fantasy worlds, among others. In 1925, Russian psychiatrist Grunya Sukhareva described a "schizoid psychopathy" in a group of children, resembling today's SzPD and ASD. About a decade later Pyotr Gannushkin also included Schizoids and Dreamers in his detailed typology of personality types.
The descriptive tradition began in 1925 with the description of observable schizoid behaviors by Ernst Kretschmer. He organized those into three groups of characteristics:
Unsociability, quietness, reservedness, seriousness and eccentricity.
Timidity, shyness with feelings, sensitivity, nervousness, excitability, fondness of nature and books.
Pliability, kindliness, honesty, indifference, silence and cold emotional attitudes.
These characteristics were the precursors of the DSM-III division of the schizoid character into three distinct personality disorders: schizotypal, avoidant and schizoid. Kretschmer himself, however, did not conceive of separating these behaviors to the point of radical isolation but considered them to be simultaneously present as varying potentials in schizoid individuals. For Kretschmer, the majority of schizoid people are not either oversensitive or cold, but they are oversensitive and cold "at the same time" in quite different relative proportions, with a tendency to move along these dimensions from one behavior to the other.
The second path, that of dynamic psychiatry, began in 1924 with observations by Eugen Bleuler, who observed that the schizoid person and schizoid pathology were not things to be set apart. Ronald Fairbairn's seminal work on the schizoid personality, from which most of what is known today about schizoid phenomena is derived, was presented in 1940. Here, Fairbairn delineated four central schizoid themes:
The need to regulate interpersonal distance as a central focus of concern.
The ability to mobilize self-preservative defenses and self-reliance.
A pervasive tension between the anxiety-laden need for attachment and the defensive need for distance that manifests in observable behavior as indifference.
An overvaluation of the inner world at the expense of the outer world.
Following Fairbairn's derivation of SzPD from a combination of derealization, depersonalization, splitting, the oral stage of making all subjects into partial objects, and intellectualization; the dynamic psychiatry tradition has continued to produce rich explorations on the schizoid character, most notably from writers Nannarello (1953), Laing (1965), Winnicott (1965), Guntrip (1969), Khan (1974), Akhtar (1987), Seinfeld (1991), Manfield (1992) and Klein (1995).
The DSM-I had the diagnosis of schizoid personality, which was defined by avoidance of close relationships, inability to express aggressive feelings, and autistic thinking (thinking which is preoccupied with one's inner experience). The DSM-II later updated the definition to include daydreaming, detachment from reality, and sensitivity. It was incorporated into the DSM-III as schizoid personality disorder to describe difficulties forming meaningful social relationships and a persistent pattern of disconnection and apathy. The diagnosis of SzPD made it to the DSM-IV and DSM-V.
Epidemiology
It remains unclear how prevalent the disorder is. It may be present in anywhere from 0.5% to 7% of the population and possibly 14% of the homeless population. Gender differences in this disorder are also unclear. Some research has suggested that this disorder may occur more frequently in men than women. SzPD is uncommon in clinical settings (about 2.2%) and occurs more commonly in males. It is rare compared with other personality disorders. Philip Manfield suggests that the "schizoid condition", which roughly includes the DSM schizoid, avoidant and schizotypal personality disorders, is represented by "as many as forty percent of all personality disorders." Manfield adds: "This huge discrepancy [from the ten percent reported by therapists for the condition] is probably largely because someone with a schizoid disorder is less likely to seek treatment than someone with other axis-II disorders." A 2008 study assessing personality and mood disorder prevalence among homeless people at New York City drop-in centers reported an SzPD rate of 65% among this sample. The study did not assess homeless people who did not show up at drop-in centers, and the rates of most other personality and mood disorders within the drop-in centers were lower than that of SzPD. The authors noted the limitations of the study, including the higher male-to-female ratio in the sample and the absence of subjects outside the support system or receiving other support (e.g., shelters) as well as the absence of subjects in geographical settings outside New York City, a large city often considered a magnet for disenfranchised people. A University of Colorado Colorado Springs study comparing personality disorders and Myers–Briggs Type Indicator types found that the disorder had a significant correlation with the Introverted (I) and Thinking (T) preferences.
Etiology
Environmental
Perfectionist and hypercritical parenting or cold, neglectful, and distant parenting contribute to the onset of SzPD. For a person with SzPD, their parents likely were intolerant of their emotional experiences. They may have been forced to repress and compartmentalize their emotions, possibly resulting in the onset of difficulties expressing and processing emotional experiences. These difficulties lead to the child feeling rejected and developing the belief that the only safe environment is one where they are alone and inexpressive. People with SzPD may also have internalized the belief that their emotions are dangerous to themselves and others due to the negative responses received from others. In their status of isolation and emotional bluntness they can be self-sufficient and safe. Childhood trauma can also contribute to feelings of emptiness in adulthood. Alcoholism in parents is associated with a heightened risk of developing SzPD.
Genetic
Sula Wolff, who did extensive research and clinical work with children and teenagers with schizoid symptoms, stated that "schizoid personality has a constitutional, probably genetic, basis." Research on heritability and this disorder is lacking. Twin studies with SzPD traits (e.g., low sociability and low warmth) suggest that these traits are inherited. Besides this indirect evidence, the direct heritability estimates of SzPD range from 50% to 59%. Earlier, less methodologically rigorous research had found the heritability rate to be 29%.
The pathophysiology of SzPD remains unclear. Genetic relationships with people who have schizophrenia spectrum disorders increase the risk of developing schizoid personality disorder. People with SzPD can have a history of schizotypy before developing the disorder. SzPD symptoms can be premorbid to schizophrenia.
Neurological
Prenatal malnutrition, premature birth, and low birth weight are all thought to play a role in the development of SzPD. SzPD is associated with reduced serotonergic and dopaminergic pathways in areas such as the frontal lobe, amygdala, and striatum. Traumatic brain injuries to the frontal lobe may also contribute to the onset of SzPD as that area of the brain controls areas such as emotion and socialization. Deficits in the right hemisphere of the brain may also be associated with SzPD. Lower levels of low-density lipoprotein cholesterol may be correlated with the presence of schizoid traits in women. Excess indices in the left hemisphere may also be related to SzPD.
Prognosis
Traits of schizoid personality disorder appear in childhood and adolescence. Children with this disorder usually have poor relationships with others, social anxiety, internal fantasies, strange behavior, and hyperactivity. These behaviors can result in teasing and bullying at the hands of others. It is common for people with SzPD to have had major depressive disorder in childhood. SzPD is associated with lower levels of achievement, a compromised quality of life, and a worse outcome of treatment. Treatment for this disorder is under-studied and poorly understood. There is no widely accepted and approved psychotherapy or medication for this disorder. It is one of the most poorly researched psychiatric disorders. Professionals may misunderstand the disorder and the client, potentially reinforcing a feeling of failure and negatively impacting their willingness to continue to commit to treatment. Clinicians tend to worry that they are incapable of properly treating the patient. It is rare for someone with this disorder to voluntarily seek treatment without a comorbid disorder or pressure from family or friends. In treatment, people with SzPD are usually disinterested and often minimize symptoms. Patients with SzPD may fear losing their independence through therapy. Many schizoid individuals will avoid making the efforts required to establish a proper relationship with the therapist. It can be difficult for them to open up or discuss their emotions in therapy. Although people with this disorder can still improve, it is unlikely they will ever experience significant joy through social interaction.
Signs and symptoms
Social isolation
SzPD is associated with a dismissive-avoidant attachment style. People with this disorder will rarely maintain close relationships and often exclusively choose to participate in solitary activities. People with schizoid personality disorder typically have no close friends or confidants, except for a close relative on occasions.
They usually prefer hobbies and activities that do not require interaction with others. People with SzPD may be averse to social situations due to difficulties deriving pleasure from physical or emotional sensations, rather than social anhedonia.
One potential motivation for avoiding social situations is that they feel that it intrudes on their freedom. Relationships can feel suffocating for people with SzPD, and they may think of them as opportunities for entrapment.
Patients with this disorder are often independent and turn to themselves as sources of validation. They tend to be the happiest when in relationships in which their partner places few emotional or intimate demands on them and does not expect phatic or social niceties. It is not necessarily people they want to avoid, but negative or positive emotional expectations, emotional intimacy, and self-disclosure.
Patients with SzPD can feel as if close emotional bonds are dangerous to themselves and others. They may have feelings of inadequacy or shame. Some people with SzPD may experience a deep desire to connect with others, yet will be terrified by the dangers inherent in doing so. Avoidance of social situations may be a method of avoiding being hurt or rejected.
Individuals with SzPD can form relationships with others based on intellectual, physical, familial, occupational, or recreational activities, as long as there is no need for emotional intimacy. Donald Winnicott explains this is because schizoid individuals "prefer to make relationships on their own terms and not in terms of the impulses of other people." Failing to attain that, they prefer isolation.
In general, friendship for schizoid individuals is usually limited to one other person, who is often also schizoid, forming what has been called a union of two eccentrics; "within it – the ecstatic cult of personality, outside it – everything is sharply rejected and despised". Their unique lifestyle can lead to social rejection and people with SzPD are at a higher risk of facing bullying or homelessness. This social rejection can reinforce their asocial behavior.
Sexuality
People with this disorder usually have little to no interest in sexual or romantic relationships. They rarely date or marry. Sex often causes individuals with SzPD to feel that their personal space is being violated, and they commonly feel that masturbation or sexual abstinence is preferable to the emotional closeness they must tolerate when having sex. Significantly broadening this picture are notable exceptions of SzPD individuals who engage in occasional or even frequent sexual activities with others. Individuals with SzPD have long been noted to have an increased rate of unconventional sexual tendencies, though if present, these are rarely acted upon. Schizoid people are often labeled asexual or present with "a lack of sexual identity". Kernberg states that this apparent lack of sexuality does not represent a lack of sexual definition but rather a combination of several strong fixations to cope with the same conflicts. People with SzPD are often able to pursue any fantasies with content on the internet while remaining completely unengaged with the outside world.
Emotions
Sensory or emotional experiences typically provide little enjoyment for people with SzPD. They rarely display strong emotions or react to anything. People with SzPD can have difficulty expressing themselves and seem to be directionless or passive. Individuals with SzPD can also experience anhedonia. They can also have difficulty understanding others' emotions and social cues. It can be hard for people with SzPD to assess the impact of their actions in social situations. People with this condition are often indifferent towards criticism or praise and can appear distant, aloof, or uncaring to others. They may avoid others and expressing themselves as a method of keeping others distant and preventing themselves from being hurt. Remaining alone and expressionless can feel safe and comfortable for people with SzPD. Expressing themselves can make them feel shame or discomfort. People with SzPD may feel inadequate and can be sensitive, although they have difficulty expressing it. Alexithymia, or difficulties understanding one's own emotions, is common amongst people with SzPD. This leads to them isolating themselves to avoid the discomfort and stimulation that emotional experiences offer. According to Guntrip, Klein, and others, people with SzPD may possess a hidden sense of superiority and lack dependence on other people's opinions. This is very different from the grandiosity seen in narcissistic personality disorder, which is described as "burdened with envy" and with a desire to destroy or put down others. Additionally, schizoid individuals do not go out of their way to achieve social validation. Unlike narcissists, schizoid people will often keep their creations private to avoid unwelcome attention or the feeling that their ideas and thoughts are being appropriated by the public. When forced to rely on others, a person with SzPD may feel panic or terror.
Feelings of unreality
Patients with SzPD often feel unreal, empty, and separate from their own emotions. They tend to perceive themselves as fundamentally different from others and can believe that they are fundamentally unlikeable. Other people often seem strange and incomprehensible to a person with SzPD. Reality can feel unenjoyable and uninteresting to people with SzPD. They have difficulty finding motivation and lack ambition. Patients with SzPD often feel as if they are "going through the motions" or that "life passes them by." Many describe feeling as if they are observing life from a distance. Aaron Beck and his colleagues report that people with SzPD seem comfortable with their aloof lifestyle and consider themselves observers, rather than participants in the world around them. But they also mention that many of their schizoid patients recognize themselves as socially deviant (or even defective) when confronted with the different lives of ordinary people – especially when they read books or see movies focusing on relationships. Even when schizoid individuals may not long for closeness, they can become weary of being "on the outside, looking in". These feelings may lead to depression, depersonalization, or derealization. If they do, schizoid people often experience feeling "like a robot" or "going through life in a dream". People with SzPD may try to avoid all physical activity in order to become nobody and disconnect from reality. This can lead to the patient spending a large quantity of time sleeping and ignoring bodily functions such as hygiene.
Internal fantasy
Although this disorder does not affect the patient's capacity to understand reality, they may engage in excessive daydreaming and introspection. Their daydreams can grow to consume most of their lives. Real life can become secondary to their fantasy, and they can have complex lives and relationships which exist entirely inside of their internal fantasy. These daydreams may constitute a defense mechanism to protect the patient from the outside world and its difficulties. Common themes in their internal fantasies are omnipotence and grandiosity. The related schizotypal personality disorder and schizophrenia are reported to have ties to creative thinking, and it is speculated that the internal fantasy aspect of SzPD may also be reflective of this thinking. Alternatively, there has been an especially large contribution of people with schizoid symptoms to science and theoretical areas of knowledge, including mathematics, physics, economics, etc. At the same time, people with SzPD are helpless at many practical activities because of their symptoms.
Suicide and self-harm
Symptoms of SzPD such as isolation and the blunted affect put people with schizoid personality disorder at a higher risk of suicide and non-suicidal self-harm. This may be because their reduced capacities for emotion prevent them from properly dealing with strife. Their solitary nature may contribute by preventing them from finding relief in relationships. Demonstrative suicides or suicide blackmail, as seen in cluster B personality disorders such as borderline, histrionic, or antisocial, are extremely rare among schizoid individuals. As in other clinical mental health settings, among suicidal inpatients, individuals with SzPD are not as well represented as some other groups. A 2011 study on suicidal inpatients at a Moscow hospital found that schizoid individuals were the least common patients, while those with cluster B personality disorders were the most common.
Low weight
A study that looked at the body mass index (BMI) of a sample of both male adolescents diagnosed with SzPD and those diagnosed with Asperger syndrome found that the BMI of all patients was significantly below normal. Clinical records indicated abnormal eating behavior by some patients. Some patients would only eat when alone and refused to eat out. Restrictive diets and fears of disease were also found. It was suggested that the anhedonia of SzPD may also affect eating, leading schizoid individuals to not enjoy it. Alternatively, it was suggested that schizoid individuals may not feel hunger as strongly as others or not respond to it, a certain withdrawal "from themselves".
Substance abuse
Very little data exists for rates of substance use disorder among people with SzPD, but existing studies suggest they are less likely to have substance abuse problems than the general population. One study found that significantly fewer boys with SzPD had alcohol problems than a control group of non-schizoid people. Another study evaluating personality disorder profiles in substance abusers found that substance abusers who showed schizoid symptoms were more likely to abuse one substance rather than many, in contrast to other personality disorders such as borderline, antisocial, or histrionic, which were more likely to abuse many. American psychotherapist Sharon Ekleberry states that the impoverished social connections experienced by people with SzPD limit their exposure to the drug culture and that they have limited inclination to learn how to do illegal drugs. Describing them as "highly resistant to influence", she additionally states that even if they could access illegal drugs, they would be disinclined to use them in public or social settings, and because they would be more likely to use alcohol or cannabis alone than for social disinhibition, they would not be particularly vulnerable to negative consequences in early use. People with SzPD are at a lower risk of substance abuse issues than people with other personality disorders. They may form relationships with their substances as a substitute for human contact or to cope with emotional issues. People with SzPD may desire psychedelic drugs more than other kinds.
Secret schizoids
Many schizoid individuals display an engaging, interactive personality, contradicting the observable characteristic emphasized by the DSM-5 and ICD-10 definitions of the schizoid personality. Guntrip (using ideas of Klein, Fairbairn, and Winnicott) classifies these individuals as "secret schizoids", who behave with socially available, interested, engaged, and involved interaction yet remain emotionally withdrawn and sequestered within the safety of the internal world. Klein distinguishes between a "classic" SzPD and a "secret" SzPD, which occur "just as often" as each other. Klein cautions one should not misidentify the schizoid person as a result of the patient's defensive, compensatory interaction with the external world. He suggests one ask the person what their subjective experience is, to detect the presence of the schizoid refusal of emotional intimacy and preference for objective fact. A 2013 study looking at personality disorders and Internet use found that being online more hours per day predicted signs of SzPD. Additionally, SzPD correlated with lower phone call use and fewer Facebook friends.
Descriptions of the schizoid personality as "hidden" behind an outward appearance of emotional engagement have been recognized since 1940, with Fairbairn's description of "schizoid exhibitionism", in which the schizoid individual can express a great deal of feeling and make what appear to be impressive social contacts yet, in reality, gives nothing and loses nothing. Because they are "playing a part", their personality is not involved. According to Fairbairn, the person disowns the part they are playing, and the schizoid individual seeks to preserve their personality intact and immune from compromise. The schizoid person's false persona is based on what those around them define as normal or good behavior, as a form of compliance. Further references to the secret schizoid come from Masud Khan, Jeffrey Seinfeld, and Philip Manfield. These scholars described secret schizoids as people who enjoy public speaking engagements but experience great difficulty during the breaks when audience members would attempt to engage them emotionally. These references expose the problems in relying on outer observable behavior for assessing the presence of personality disorders in certain individuals.
Comorbid disorders
Agoraphobia
Avoidant personality disorder
Antisocial personality disorder
Borderline personality disorder
Post-traumatic stress disorder
Major depressive disorder
Generalized anxiety disorder
Panic disorder
Paranoid personality disorder
Social anxiety disorder
Schizotypal personality disorder
Obsessive–compulsive disorder
Autism spectrum disorder
Several studies have reported an overlap or comorbidity with autism spectrum disorder and Asperger syndrome. Asperger syndrome had traditionally been called "schizoid disorder of childhood", and Eugen Bleuler coined both the terms "autism" and "schizoid" to describe withdrawal to an internal fantasy, against which any influence from outside becomes an intolerable disturbance. In a 2012 study of a sample of 54 young adults with Asperger syndrome, it was found that 26% of them also met the criteria for SzPD, the highest comorbidity out of any personality disorder in the sample (the other comorbidities were 19% for obsessive–compulsive personality disorder, 13% for avoidant personality disorder and one female with schizotypal personality disorder). Additionally, twice as many men with Asperger syndrome met the criteria for SzPD than women. While 41% of the whole sample were unemployed with no occupation, this rose to 62% for the Asperger's and SzPD comorbid group. Tantam suggested that Asperger syndrome may confer an increased risk of developing SzPD. A 2019 study found that 54% of a group of males aged 11 to 25 with Asperger syndrome showed significant SzPD traits, with 6% meeting full diagnostic criteria for SzPD, compared to 0% of a control group.
In the 2012 study, it was noted that the DSM may complicate diagnosis by requiring the exclusion of a pervasive developmental disorder (PDD) before establishing a diagnosis of SzPD. The study found that social interaction impairments, stereotyped behaviors, and specific interests were more severe in the individuals with Asperger syndrome also fulfilling SzPD criteria, against the notion that social interaction skills are unimpaired in SzPD. The authors believe that a substantial subgroup of people with autism spectrum disorder or PDD have clear "schizoid traits" and correspond largely to the "loners" in Lorna Wing's classification The autism spectrum (Lancet 1997), described by Sula Wolff. The authors of the 2019 study hypothesized that it is extremely likely that historic cohorts of adults diagnosed with SzPD either also had childhood-onset autistic syndromes or were misdiagnosed. They stressed that further research to clarify overlap and distinctions between these two syndromes was strongly warranted, especially given that high-functioning autism spectrum disorders are now recognized in around 1% of the population.
Treatment
Medication
There are no effective medications for schizoid personality disorder. However, certain medications may reduce the symptoms of SzPD and treat co-occurring mental disorders. Since the symptoms of SzPD mirror the negative symptoms of schizophrenia, antipsychotics have been suggested as a potentially effective medication for SzPD. Originally, low doses of atypical antipsychotics like risperidone or olanzapine were used to alleviate social deficits and blunted affect. However, a 2012 review concluded that atypical antipsychotics were ineffective for treating personality disorders. Antidepressants, SSRIs, anxiolitics, bupropion, modafinil, benzodiazepines, and biofeedback may also be effective treatments.
Psychotherapy
Treatment for this disorder uses a combination of cognitive-behavioral therapy and psychodynamic psychotherapy. These techniques can be used to help patients identify their defense mechanisms and change them. Therapists attempt to establish healthy relationships with their clients, helping to combat their internalized belief that relationships are harmful and unhelpful. Relationships with a therapist can seem terrifying and intrusive to a person with SzPD. They may feel as if they need to alter or hide their feelings to meet the therapist's demands or expectations. To combat this, therapists try to gradually increase their patient's emotional expression. Expressing too much too early can lead to their ending therapy. Treatment must be person centered, with clients feeling understood and well regarded. This can allow them to connect with and understand their emotions. When people with SzPD do not have their feelings validated, this will confirm their belief that expressing themselves is dangerous. Therapists attempt to avoid intruding on their patients' lives or restricting their freedoms, so as to prevent them from feeling as if therapy is intolerable. Because of this, therapy is usually less structured than treatment programs for other disorders. Patients may benefit from long-term treatment lasting several years. Inpatient care may be effective for treating SzPD and other Cluster A disorders.
Controversy
The original concept of the schizoid character developed by Ernst Kretschmer in the 1920s comprised a mix of avoidant, schizotypal, and schizoid traits. It was not until 1980 and the work of Theodore Millon that led to splitting this concept into three personality disorders (now schizoid, schizotypal, and avoidant). This caused debate about whether this was accurate or if these traits were different expressions of a single personality disorder. It has also been argued due to the poor consistency and efficiency of diagnosis due to overlapping traits that SzPD should be removed altogether from the DSM. A 2012 article suggested that two different disorders may better represent SzPD: one affect-constricted disorder (belonging to schizotypal PD) and a seclusive disorder (belonging to avoidant PD). They called for the replacement of the SzPD category from future editions of the DSM with a dimensional model which would allow for the description of schizoid traits on an individual basis.
Some critics such as Nancy McWilliams of Rutgers University and Panagiotis Parpottas of European University Cyprus argue that the definition of SzPD is flawed due to cultural bias and that it does not constitute a mental disorder but simply an avoidant attachment style requiring a more distant emotional proximity. If that is true, then many of the more problematic reactions these individuals show in social situations may be partly accounted for by the judgments commonly imposed on people with this style.
Similarly, John Oldham, using a dimensional approach, thinks that most people with schizoid character features do not have a full-blown personality disorder. Impairment is mandatory for any behavior to be diagnosed as a personality disorder.
Diagnosis
Guntrip criteria
Ralph Klein, Clinical Director of the Masterson Institute, delineates the following nine characteristics of the schizoid personality as described by Harry Guntrip:
Introversion
Withdrawnness
Narcissism
Self-sufficiency
A sense of superiority
Loss of affect
Loneliness
Depersonalization
Regression
The description of Guntrip's nine characteristics should clarify some differences between the traditional DSM portrait of SzPD and the traditional informed object relations view. All nine characteristics are consistent. Most, if not all, must be present to diagnose a schizoid disorder.
Millon's subtypes
Theodore Millon restricted the term "schizoid" to those personalities who lack the capacity to form social relationships. He characterizes their way of thinking as being vague and void of thoughts and as sometimes having a "defective perceptual scanning". Because they often do not perceive cues that trigger affective responses, they experience fewer emotional reactions.
For Millon, SzPD is distinguished from other personality disorders in that it is "the personality disorder that lacks a personality." He criticizes that this may be due to the current diagnostic criteria: They describe SzPD only by an absence of certain traits, which results in a "deficit syndrome" or "vacuum". Instead of delineating the presence of something, they mention solely what is lacking. Therefore, it is hard to describe and research such a concept.
He identified four subtypes of SzPD. Any schizoid individual may exhibit none or one of the following:
Akhtar's profile
American psychoanalyst Salman Akhtar provided a comprehensive phenomenological profile of SzPD in which classic and contemporary descriptive views are synthesized with psychoanalytic observations. This profile is summarized in the table reproduced below that lists clinical features that involve six areas of psychosocial functioning and are organized by "overt" and "covert" manifestations.
"Overt" and "covert" are intended to denote seemingly contradictory aspects that may both simultaneously be present in an individual. These designations do not necessarily imply their conscious or unconscious existence. The covert characteristics are by definition difficult to discern and not immediately apparent. Additionally, the lack of data on the frequency of many of the features makes their relative diagnostic weight difficult to distinguish at this time. However, Akhtar states that his profile has several advantages over the DSM in terms of maintaining historical continuity of the use of the word schizoid, valuing depth and complexity over descriptive oversimplification and helping provide a more meaningful differential diagnosis of SzPD from other personality disorders.
Differential diagnosis
| Biology and health sciences | Mental disorders | Health |
229104 | https://en.wikipedia.org/wiki/Matter%20wave | Matter wave | Matter waves are a central part of the theory of quantum mechanics, being half of wave–particle duality. At all scales where measurements have been practical, matter exhibits wave-like behavior. For example, a beam of electrons can be diffracted just like a beam of light or a water wave.
The concept that matter behaves like a wave was proposed by French physicist Louis de Broglie () in 1924, and so matter waves are also known as de Broglie waves.
The de Broglie wavelength is the wavelength, , associated with a particle with momentum through the Planck constant, :
Wave-like behavior of matter has been experimentally demonstrated, first for electrons in 1927 and for other elementary particles, neutral atoms and molecules in the years since.
Matter waves have more complex velocity relations than solid objects and they also differ from electromagnetic waves (light). Collective matter waves are used to model phenomena in solid state physics; standing matter waves are used in molecular chemistry.
Matter wave concepts are widely used in the study of materials where different wavelength and interaction characteristics of electrons, neutrons, and atoms are leveraged for advanced microscopy and diffraction technologies.
History
Background
At the end of the 19th century, light was thought to consist of waves of electromagnetic fields which propagated according to Maxwell's equations, while matter was thought to consist of localized particles (see history of wave and particle duality). In 1900, this division was questioned when, investigating the theory of black-body radiation, Max Planck proposed that the thermal energy of oscillating atoms is divided into discrete portions, or quanta. Extending Planck's investigation in several ways, including its connection with the photoelectric effect, Albert Einstein proposed in 1905 that light is also propagated and absorbed in quanta, now called photons. These quanta would have an energy given by the Planck–Einstein relation:
and a momentum vector
where (lowercase Greek letter nu) and (lowercase Greek letter lambda) denote the frequency and wavelength of the light, the speed of light, and the Planck constant. In the modern convention, frequency is symbolized by as is done in the rest of this article. Einstein's postulate was verified experimentally by K. T. Compton and O. W. Richardson and by A. L. Hughes in 1912 then more carefully including a measurement of the Planck constant in 1916 by Robert Millikan.
De Broglie hypothesis
De Broglie, in his 1924 PhD thesis, proposed that just as light has both wave-like and particle-like properties, electrons also have wave-like properties.
His thesis started from the hypothesis, "that to each portion of energy with a proper mass one may associate a periodic phenomenon of the frequency , such that one finds: . The frequency is to be measured, of course, in the rest frame of the energy packet. This hypothesis is the basis of our theory." (This frequency is also known as Compton frequency.)
To find the wavelength equivalent to a moving body, de Broglie set the total energy from special relativity for that body equal to :
(Modern physics no longer uses this form of the total energy; the energy–momentum relation has proven more useful.) De Broglie identified the velocity of the particle, , with the wave group velocity in free space:
(The modern definition of group velocity uses angular frequency and wave number ). By applying the differentials to the energy equation and identifying the relativistic momentum:
then integrating, de Broglie arrived at his formula for the relationship between the wavelength, , associated with an electron and the modulus of its momentum, , through the Planck constant, :
Schrödinger's (matter) wave equation
Following up on de Broglie's ideas, physicist Peter Debye made an offhand comment that if particles behaved as waves, they should satisfy some sort of wave equation. Inspired by Debye's remark, Erwin Schrödinger decided to find a proper three-dimensional wave equation for the electron. He was guided by William Rowan Hamilton's analogy between mechanics and optics (see Hamilton's optico-mechanical analogy), encoded in the observation that the zero-wavelength limit of optics resembles a mechanical system – the trajectories of light rays become sharp tracks that obey Fermat's principle, an analog of the principle of least action.
In 1926, Schrödinger published the wave equation that now bears his name – the matter wave analogue of Maxwell's equations – and used it to derive the energy spectrum of hydrogen. Frequencies of solutions of the non-relativistic Schrödinger equation differ from de Broglie waves by the Compton frequency since the energy corresponding to the rest mass of a particle is not part of the non-relativistic Schrödinger equation. The Schrödinger equation describes the time evolution of a wavefunction, a function that assigns a complex number to each point in space. Schrödinger tried to interpret the modulus squared of the wavefunction as a charge density. This approach was, however, unsuccessful. Max Born proposed that the modulus squared of the wavefunction is instead a probability density, a successful proposal now known as the Born rule.
The following year, 1927, C. G. Darwin (grandson of the famous biologist) explored Schrödinger's equation in several idealized scenarios. For an unbound electron in free space he worked out the propagation of the wave, assuming an initial Gaussian wave packet. Darwin showed that at time later the position of the packet traveling at velocity would be
where is the uncertainty in the initial position. This position uncertainty creates uncertainty in velocity (the extra second term in the square root) consistent with Heisenberg's uncertainty relation The wave packet spreads out as show in the figure.
Experimental confirmation
In 1927, matter waves were first experimentally confirmed to occur in George Paget Thomson and Alexander Reid's diffraction experiment and the Davisson–Germer experiment, both for electrons.
The de Broglie hypothesis and the existence of matter waves has been confirmed for other elementary particles, neutral atoms and even molecules have been shown to be wave-like.
The first electron wave interference patterns directly demonstrating wave–particle duality used electron biprisms (essentially a wire placed in an electron microscope) and measured single electrons building up the diffraction pattern.
Recently, a close copy of the famous double-slit experiment using electrons through physical apertures gave the movie shown.
Electrons
In 1927 at Bell Labs, Clinton Davisson and Lester Germer fired slow-moving electrons at a crystalline nickel target. The diffracted electron intensity was measured, and was determined to have a similar angular dependence to diffraction patterns predicted by Bragg for x-rays. At the same time George Paget Thomson and Alexander Reid at the University of Aberdeen were independently firing electrons at thin celluloid foils and later metal films, observing rings which can be similarly interpreted. (Alexander Reid, who was Thomson's graduate student, performed the first experiments but he died soon after in a motorcycle accident and is rarely mentioned.) Before the acceptance of the de Broglie hypothesis, diffraction was a property that was thought to be exhibited only by waves. Therefore, the presence of any diffraction effects by matter demonstrated the wave-like nature of matter. The matter wave interpretation was placed onto a solid foundation in 1928 by Hans Bethe, who solved the Schrödinger equation, showing how this could explain the experimental results. His approach is similar to what is used in modern electron diffraction approaches.
This was a pivotal result in the development of quantum mechanics. Just as the photoelectric effect demonstrated the particle nature of light, these experiments showed the wave nature of matter.
Neutrons
Neutrons, produced in nuclear reactors with kinetic energy of around , thermalize to around as they scatter from light atoms. The resulting de Broglie wavelength (around ) matches interatomic spacing and neutrons scatter strongly from hydrogen atoms. Consequently, neutron matter waves are used in crystallography, especially for biological materials. Neutrons were discovered in the early 1930s, and their diffraction was observed in 1936. In 1944, Ernest O. Wollan, with a background in X-ray scattering from his PhD work under Arthur Compton, recognized the potential for applying thermal neutrons from the newly operational X-10 nuclear reactor to crystallography. Joined by Clifford G. Shull, they developed neutron diffraction throughout the 1940s.
In the 1970s, a neutron interferometer demonstrated the action of gravity in relation to wave–particle duality. The double-slit experiment was performed using neutrons in 1988.
Atoms
Interference of atom matter waves was first observed by Immanuel Estermann and Otto Stern in 1930, when a Na beam was diffracted off a surface of NaCl. The short de Broglie wavelength of atoms prevented progress for many years until two technological breakthroughs revived interest: microlithography allowing precise small devices and laser cooling allowing atoms to be slowed, increasing their de Broglie wavelength. The double-slit experiment on atoms was performed in 1991.
Advances in laser cooling allowed cooling of neutral atoms down to nanokelvin temperatures. At these temperatures, the de Broglie wavelengths come into the micrometre range. Using Bragg diffraction of atoms and a Ramsey interferometry technique, the de Broglie wavelength of cold sodium atoms was explicitly measured and found to be consistent with the temperature measured by a different method.
Molecules
Recent experiments confirm the relations for molecules and even macromolecules that otherwise might be supposed too large to undergo quantum mechanical effects. In 1999, a research team in Vienna demonstrated diffraction for molecules as large as fullerenes. The researchers calculated a de Broglie wavelength of the most probable C60 velocity as .
More recent experiments prove the quantum nature of molecules made of 810 atoms and with a mass of . As of 2019, this has been pushed to molecules of .
In these experiments the build-up of such interference patterns could be recorded in real time and with single molecule sensitivity.
Large molecules are already so complex that they give experimental access to some aspects of the quantum-classical interface, i.e., to certain decoherence mechanisms.
Others
Matter wave was detected in van der Waals molecules, rho mesons, Bose-Einstein condensate.
Traveling matter waves
Waves have more complicated concepts for velocity than solid objects.
The simplest approach is to focus on the description in terms of plane matter waves for a free particle, that is a wave function described by
where is a position in real space, is the wave vector in units of inverse meters, is the angular frequency with units of inverse time and is time. (Here the physics definition for the wave vector is used, which is times the wave vector used in crystallography, see wavevector.) The de Broglie equations relate the wavelength to the modulus of the momentum , and frequency to the total energy of a free particle as written above:
where is the Planck constant. The equations can also be written as
Here, is the reduced Planck constant. The second equation is also referred to as the Planck–Einstein relation.
Group velocity
In the de Broglie hypothesis, the velocity of a particle equals the group velocity of the matter wave.
In isotropic media or a vacuum the group velocity of a wave is defined by:
The relationship between the angular frequency and wavevector is called the dispersion relationship. For the non-relativistic case this is:
where is the rest mass. Applying the derivative gives the (non-relativistic) matter wave group velocity:
For comparison, the group velocity of light, with a dispersion , is the speed of light .
As an alternative, using the relativistic dispersion relationship for matter waves
then
This relativistic form relates to the phase velocity as discussed below.
For non-isotropic media we use the Energy–momentum form instead:
But (see below), since the phase velocity is , then
where is the velocity of the center of mass of the particle, identical to the group velocity.
Phase velocity
The phase velocity in isotropic media is defined as:
Using the relativistic group velocity above:
This shows that as reported by R.W. Ditchburn in 1948 and J. L. Synge in 1952. Electromagnetic waves also obey , as both and . Since for matter waves, , it follows that , but only the group velocity carries information. The superluminal phase velocity therefore does not violate special relativity, as it does not carry information.
For non-isotropic media, then
Using the relativistic relations for energy and momentum yields
The variable can either be interpreted as the speed of the particle or the group velocity of the corresponding matter wave—the two are the same. Since the particle speed for any particle that has nonzero mass (according to special relativity), the phase velocity of matter waves always exceeds c, i.e.,
which approaches c when the particle speed is relativistic. The superluminal phase velocity does not violate special relativity, similar to the case above for non-isotropic media. See the article on Dispersion (optics) for further details.
Special relativity
Using two formulas from special relativity, one for the relativistic mass energy and one for the relativistic momentum
allows the equations for de Broglie wavelength and frequency to be written as
where is the velocity, the Lorentz factor, and the speed of light in vacuum. This shows that as the velocity of a particle approaches zero (rest) the de Broglie wavelength approaches infinity.
Four-vectors
Using four-vectors, the de Broglie relations form a single equation:
which is frame-independent.
Likewise, the relation between group/particle velocity and phase velocity is given in frame-independent form by:
where
Four-momentum
Four-wavevector
Four-velocity
General matter waves
The preceding sections refer specifically to free particles for which the wavefunctions are plane waves. There are significant numbers of other matter waves, which can be broadly split into three classes: single-particle matter waves, collective matter waves and standing waves.
Single-particle matter waves
The more general description of matter waves corresponding to a single particle type (e.g. a single electron or neutron only) would have a form similar to
where now there is an additional spatial term in the front, and the energy has been written more generally as a function of the wave vector. The various terms given before still apply, although the energy is no longer always proportional to the wave vector squared. A common approach is to define an effective mass which in general is a tensor given by
so that in the simple case where all directions are the same the form is similar to that of a free wave above.In general the group velocity would be replaced by the probability current
where is the del or gradient operator. The momentum would then be described using the kinetic momentum operator,
The wavelength is still described as the inverse of the modulus of the wavevector, although measurement is more complex. There are many cases where this approach is used to describe single-particle matter waves:
Bloch wave, which form the basis of much of band structure as described in Ashcroft and Mermin, and are also used to describe the diffraction of high-energy electrons by solids.
Waves with angular momentum such as electron vortex beams.
Evanescent waves, where the component of the wavevector in one direction is complex. These are common when matter waves are being reflected, particularly for grazing-incidence diffraction.
Collective matter waves
Other classes of matter waves involve more than one particle, so are called collective waves and are often quasiparticles. Many of these occur in solids – see Ashcroft and Mermin. Examples include:
In solids, an electron quasiparticle is an electron where interactions with other electrons in the solid have been included. An electron quasiparticle has the same charge and spin as a "normal" (elementary particle) electron and, like a normal electron, it is a fermion. However, its effective mass can differ substantially from that of a normal electron. Its electric field is also modified, as a result of electric field screening.
A hole is a quasiparticle which can be thought of as a vacancy of an electron in a state; it is most commonly used in the context of empty states in the valence band of a semiconductor. A hole has the opposite charge of an electron.
A polaron is a quasiparticle where an electron interacts with the polarization of nearby atoms.
An exciton is an electron and hole pair which are bound together.
A Cooper pair is two electrons bound together so they behave as a single matter wave.
Standing matter waves
The third class are matter waves which have a wavevector, a wavelength and vary with time, but have a zero group velocity or probability flux. The simplest of these, similar to the notation above would be
These occur as part of the particle in a box, and other cases such as in a ring. This can, and arguably should be, extended to many other cases. For instance, in early work de Broglie used the concept that an electron matter wave must be continuous in a ring to connect to the Bohr–Sommerfeld condition in the early approaches to quantum mechanics. In that sense atomic orbitals around atoms, and also molecular orbitals are electron matter waves.
Matter waves vs. electromagnetic waves (light)
Schrödinger applied Hamilton's optico-mechanical analogy to develop his wave mechanics for subatomic particles Consequently, wave solutions to the Schrödinger equation share many properties with results of light wave optics. In particular, Kirchhoff's diffraction formula works well for electron optics and for atomic optics. The approximation works well as long as the electric fields change more slowly than the de Broglie wavelength. Macroscopic apparatus fulfill this condition; slow electrons moving in solids do not.
Beyond the equations of motion, other aspects of matter wave optics differ from the corresponding light optics cases.
Sensitivity of matter waves to environmental condition.
Many examples of electromagnetic (light) diffraction occur in air under many environmental conditions. Obviously visible light interacts weakly with air molecules. By contrast, strongly interacting particles like slow electrons and molecules require vacuum: the matter wave properties rapidly fade when they are exposed to even low pressures of gas. With special apparatus, high velocity electrons can be used to study liquids and gases. Neutrons, an important exception, interact primarily by collisions with nuclei, and thus travel several hundred feet in air.
Dispersion. Light waves of all frequencies travel at the same speed of light while matter wave velocity varies strongly with frequency. The relationship between frequency (proportional to energy) and wavenumber or velocity (proportional to momentum) is called a dispersion relation. Light waves in a vacuum have linear dispersion relation between frequency: . For matter waves the relation is non-linear:
This non-relativistic matter wave dispersion relation says the frequency in vacuum varies with wavenumber () in two parts: a constant part due to the de Broglie frequency of the rest mass () and a quadratic part due to kinetic energy. The quadratic term causes rapid spreading of wave packets of matter waves.
Coherence The visibility of diffraction features using an optical theory approach depends on the beam coherence, which at the quantum level is equivalent to a density matrix approach. As with light, transverse coherence (across the direction of propagation) can be increased by collimation. Electron optical systems use stabilized high voltage to give a narrow energy spread in combination with collimating (parallelizing) lenses and pointed filament sources to achieve good coherence. Because light at all frequencies travels the same velocity, longitudinal and temporal coherence are linked; in matter waves these are independent. For example, for atoms, velocity (energy) selection controls longitudinal coherence and pulsing or chopping controls temporal coherence.
Optically shaped matter waves
Optical manipulation of matter plays a critical role in matter wave optics: "Light waves can act as refractive, reflective, and absorptive structures for matter waves, just as glass interacts with light waves." Laser light momentum transfer can cool matter particles and alter the internal excitation state of atoms.
Multi-particle experiments
While single-particle free-space optical and matter wave equations are identical, multiparticle systems like coincidence experiments are not.
Applications of matter waves
The following subsections provide links to pages describing applications of matter waves as probes of materials or of fundamental quantum properties. In most cases these involve some method of producing travelling matter waves which initially have the simple form , then using these to probe materials.
As shown in the table below, matter wave mass ranges over 6 orders of magnitude and energy over 9 orders but the wavelengths are all in the picometre range, comparable to atomic spacings. (Atomic diameters range from 62 to 520 pm, and the typical length of a carbon–carbon single bond is 154 pm.) Reaching longer wavelengths requires special techniques like laser cooling to reach lower energies; shorter wavelengths make diffraction effects more difficult to discern. Therefore, many applications focus on material structures, in parallel with applications of electromagnetic waves, especially X-rays. Unlike light, matter wave particles may have mass, electric charge, magnetic moments, and internal structure, presenting new challenges and opportunities.
Electrons
Electron diffraction patterns emerge when energetic electrons reflect or penetrate ordered solids; analysis of the patterns leads to models of the atomic arrangement in the solids.
They are used for imaging from the micron to atomic scale using electron microscopes, in transmission, using scanning, and for surfaces at low energies.
The measurements of the energy they lose in electron energy loss spectroscopy provides information about the chemistry and electronic structure of materials. Beams of electrons also lead to characteristic X-rays in energy dispersive spectroscopy which can produce information about chemical content at the nanoscale.
Quantum tunneling explains how electrons escape from metals in an electrostatic field at energies less than classical predictions allow: the matter wave penetrates of the work function barrier in the metal.
Scanning tunneling microscope leverages quantum tunneling to image the top atomic layer of solid surfaces.
Electron holography, the electron matter wave analog of optical holography, probes the electric and magnetic fields in thin films.
Neutrons
Neutron diffraction complements x-ray diffraction through the different scattering cross sections and sensitivity to magnetism.
Small-angle neutron scattering provides way to obtain structure of disordered systems that is sensitivity to light elements, isotopes and magnetic moments.
Neutron reflectometry is a neutron diffraction technique for measuring the structure of thin films.
Neutral atoms
Atom interferometers, similar to optical interferometers, measure the difference in phase between atomic matter waves along different paths.
Atom optics mimic many light optic devices, including mirrors, atom focusing zone plates.
Scanning helium microscopy uses He atom waves to image solid structures non-destructively.
Quantum reflection uses matter wave behavior to explain grazing angle atomic reflection, the basis of some atomic mirrors.
Quantum decoherence measurements rely on Rb atom wave interference.
Molecules
Quantum superposition revealed by interference of matter waves from large molecules probes the limits of wave–particle duality and quantum macroscopicity.
Matter-wave interfererometers generate nanostructures on molecular beams that can be read with nanometer accuracy and therefore be used for highly sensitive force measurements, from which one can deduce a plethora or properties of individualized complex molecules.
| Physical sciences | Quantum mechanics | Physics |
229116 | https://en.wikipedia.org/wiki/Veterinarian | Veterinarian | A veterinarian (vet) or veterinary surgeon is a medical professional who practices veterinary medicine. They manage a wide range of health conditions and injuries in non-human animals. Along with this, veterinarians also play a role in animal reproduction, health management, conservation, husbandry and breeding and preventive medicine like nutrition, vaccination and parasitic control as well as biosecurity and zoonotic disease surveillance and prevention.
Description
In many countries, the local nomenclature for a veterinarian is a regulated and protected term, meaning that members of the public without the prerequisite qualifications and/or license are not able to use the title. This title is selective in order to produce the most knowledgeable veterinarians that pass these qualifications. In many cases, the activities that may be undertaken by a veterinarian (such as treatment of illness or surgery in animals) are restricted only to those professionals who are registered as a veterinarian. For instance, in the United Kingdom, as in other jurisdictions, animal treatment may only be performed by registered veterinarians (with a few designated exceptions, such as paraveterinary workers), and it is illegal for any person who is not registered to call themselves a veterinarian, prescribe any drugs, or perform treatment.
Most veterinarians work in a clinical setting or bricks and mortar practice, treating animals directly. Other vets work as mobile vets offering veterinary services and treating patients in their clients home. Veterinarians may be involved in general practice, treating animals of all types; they may be specialized in a specific group of animals such as companion animals, livestock, zoo animals or equines; or may specialize in a narrow medical discipline such as surgery, dermatology or internal medicine. As with other healthcare professionals, veterinarians face ethical decisions about the care of their patients. Current debates within the profession include the ethics of certain procedures believed to be purely cosmetic or unnecessary for behavioral issues, such as declawing of cats, docking of tails, cropping of ears and debarking on dogs.
Etymology and nomenclature
The word "veterinary" comes from the Latin meaning "working animals". "Veterinarian" was first used in print by Thomas Browne in 1646. Although "vet" is commonly used as an abbreviation in all English-speaking countries, the occupation is formally referred to as a veterinary surgeon in the United Kingdom and Ireland and now as a veterinarian in most of the rest of the English-speaking world.
History
Ancient Indian sage and veterinarian Shalihotra (mythological estimate c. 2350 BCE), the son of a sage, Hayagosha, is considered the founder of veterinary sciences.
The first veterinary college was founded in Lyon, France, in 1762 by Claude Bourgelat. According to Lupton, after observing the devastation being caused by cattle plague to the French herds, Bourgelat devoted his time to seeking out a remedy. This resulted in his founding a veterinary college in Lyon in 1761, from which establishment he dispatched students to combat the disease; in a short time, the plague was stayed and the health of stock restored, through the assistance rendered to agriculture by veterinary science and art.
The Odiham Agricultural Society was founded in 1783 in England to promote agriculture and industry, and played an important role in the foundation of the veterinary profession in Britain. A 1785 Society meeting resolved to "promote the study of Farriery upon rational scientific principles."
The professionalization of the veterinary trade was finally achieved in 1790, through the campaigning of Granville Penn, who persuaded the Frenchman Charles Vial de Sainbel to accept the professorship of the newly established Veterinary College in London. The Royal College of Veterinary Surgeons was established by royal charter in 1844.
Veterinary science came of age in the late 19th century, with notable contributions from Sir John McFadyean, credited by many as having been the founder of modern Veterinary research.
Roles and responsibilities
Veterinarians treat disease, disorder or injury in animals, which includes diagnosis, treatment and aftercare. The scope of practice, specialty and experience of the individual veterinarian will dictate exactly what interventions they perform, but most will perform surgery (of differing complexity).
Unlike in human medicine, veterinarians must rely primarily on clinical signs, as animals are unable to vocalize symptoms as a human would. In some cases, owners may be able to provide a medical history and the veterinarian can combine this information along with observations, and the results of pertinent diagnostic tests such as radiography, CT scans, MRI, blood tests, urinalysis and others.
Veterinarians must consider the appropriateness of euthanasia ("putting to sleep") if a condition is likely to leave the animal in pain or with a poor quality of life, or if treatment of a condition is likely to cause more harm to the patient than good, or if the patient is unlikely to survive any treatment regimen. Additionally, there are scenarios where euthanasia is considered due to the constraints of the client's finances.
As with human medicine, much veterinary work is concerned with prophylactic treatment, in order to prevent problems occurring in the future. Common interventions include vaccination against common animal illnesses, such as distemper or rabies, and dental prophylaxis to prevent or inhibit dental disease. This may also involve owner education so as to avoid future medical or behavioral issues.
Additionally, veterinarians can play important roles in public health and the prevention of zoonoses.
Employment
The majority of veterinarians are employed in private practice treating animals (75% of vets in the United States, according to the American Veterinary Medical Association).
Small animal veterinarians typically work in veterinary clinics, veterinary hospitals, or both. Large animal veterinarians often spend more time travelling to see their patients at the primary facilities which house them, such as zoos or farms.
Other employers include charities treating animals, colleges of veterinary medicine, research laboratories, animal food companies, and pharmaceutical companies. In many countries, the government may also be a major employer of veterinarians, such as the United States Department of Agriculture or the Animal and Plant Health Agency in the United Kingdom. State and local governments also employ veterinarians.
The COVID-19 pandemic has created a greater demand for veterinary services. Many people are home with extra time on their hands, and adoption agencies and animals shelters have seen a surge in pet purchases as a result. The American Veterinary Medical Association has provided COVID-19 resources for veterinarians on prevention measures, animal testing, and wellbeing.
Focus of practice
Veterinarians and their practices may be specialized in certain areas of veterinary medicine. Areas of focus include:
Exotic animal veterinarian – Specializes in treating animals other than common pets and livestock. Includes reptiles, exotic birds such as parrots and cockatoos, and small mammals such as ferrets, rabbits, and chinchillas.
Conservation medicine – The study of the relationship between animal and human health and environmental information.
Small animal practice – Usually dogs, cats, and other companion animals/household pets such as hamsters and gerbils. Some practices are canine-only or feline-only practices.
Laboratory animal practice – Some veterinarians work in a university or industrial laboratory and are responsible for the care and treatment of laboratory animals of any species (often involving bovines, porcine species, felines, canines, rodents, and even exotic animals). Their responsibility is not only for the health and well-being of the animals, but also for enforcing humane and ethical treatment of the animals in the facility.
Large animal practice – Usually referring to veterinarians that work with, variously, livestock and other large farm animals, as well as equine species and large reptiles.
Equine medicine – Some veterinarians are specialists in equine medicine. Horses are different in anatomy, physiology, pathology, pharmacology, and husbandry to other domestic species. Specialization in equine veterinary practice is something that is normally developed after qualification, even if students do have some interest before graduation.
Food supply medicine – Some veterinarians deal exclusively or primarily with animals raised for food (such as meat, milk, and eggs). Livestock practitioners may deal with ovine (sheep), bovine (cattle) and porcine (swine) species; such veterinarians deal with management of herds, nutrition, reproduction, and minor field surgery. Dairy medicine practice focuses on dairy animals. Poultry medicine practice focuses on the health of flocks of poultry; the field often involves extensive training in pathology, epidemiology, and nutrition of birds. The veterinarian treats the flock and not the individual animals.
Food safety practice – Veterinarians are employed by both the food industry and government agencies to advise on and monitor the handling, preparation, and storage of food in ways that prevent foodborne illness.
Wildlife medicine – A relatively recent branch of veterinary medicine, focusing on wildlife. Wildlife medicine veterinarians may work with zoologists and conservation medicine practitioners and may also be called out to treat marine species such as sea otters, dolphins, or whales after a natural disaster or oil spill.
Aquatic medicine – mostly refers to veterinary care of fish in aquaculture (like salmon, cod, among other species), but can also include care of aquatic mammals. For certain countries with high economic income from aquaculture, this is an important part of the veterinary field (like Norway, Chile). Other countries (particularly those which are landlocked), might have little or no emphasis on aquatic medicine.
Dentistry – Many practices are incorporating dentistry into their daily medical services. Veterinary dentistry can extend the life of the patient by preventing oral disease and keeping the teeth and gums of the patient in healthy condition.
Veterinary specialties
Veterinary specialists are in the minority compared to general practice veterinarians, and tend to be based at points of referral, such as veterinary schools or larger animal hospitals. Unlike human medicine, veterinary specialties often combine both the surgical and medical aspects of a biological system.
Veterinary specialties are accredited in North America by the AVMA through the American Board of Veterinary Specialties, in Europe by the European Board of Veterinary Specialisation and in Australia and New Zealand by the Australasian Veterinary Boards Council. While some veterinarians may have areas of interest outside of recognized specialties, they are not legally specialists.
Specialties can cover general topics such as anesthesiology, dentistry, and surgery, as well as organ system focus such as cardiology or dermatology. A full list can be seen at veterinary specialties.
Mobile practice
Many veterinarians, especially in large animal practice, offer house calls and farm calls through a mobile practice. The start-up and operating costs of a mobile practice are typically lower than those of a traditional brick and mortar hospital, which can cost millions of dollars or more for equipment and surgical supplies. Costs associated with mobile units can range from as low as $5,000 for a utility box in an SUV to around $250,000 for a fully equipped custom built chassis. The potential advantages to the client are not having to transport the animal, lower stress for the animal, a lower risk of disease transmission from other animals, and convenience. A 2015 study published in the Journal of American Veterinary Medical Association proved that blood pressure readings, pulse rates and body temperature rates were increased by 11–16% when those readings were done in the clinic versus in the home. However, mobile practices often lack the facilities and equipment to provide advanced care, surgery, or hospitalization. Some mobile practices maintain a relationship with a traditional hospital for referral of cases needing more comprehensive care.
Salary
The last AVMA Report on Veterinary Compensation, published in 2018, indicated private practice associate veterinarians who had board certification earned a mean of $187,000. A veterinarian's salary can easily exceed $300,000 depending on the specialty. The median starting salary for new veterinary graduates without specialization in 2018 was $103,800 in the United States according to the Bureau of Labor Statistics, while the lowest paid earned less than $89,540 annually. States and districts with the highest mean salary are California ($398,340), Michigan ($325,100), Illinois ($324,870), New York ($322,500), and Hawaii ($221,150). Veterinarians who own their own clinics are typically paid a much higher salary. The average owner payout is $400,000 for every $1,000,000 of clinic income. In 2021 there were practices sold with $8–10,000,000 in yearly revenue with the owners drawing salaries of several million dollars. Over 90% of practice owners do not regret purchasing or starting their own practice, according to a 2020 survey of clinic owners.
Education and regulation
In order to practice, vets must complete an appropriate degree in veterinary medicine, and in most cases must also be registered with the relevant governing body for their jurisdiction.
Veterinary science degrees
Degrees in veterinary medicine culminate in the award of a veterinary science degree, although the title varies by region. For instance, in North America, graduates will receive a Doctor of Veterinary Medicine (Doctor of Veterinary Medicine or Veterinariae Medicinae Doctoris; DVM or VMD), whereas in the United Kingdom, Australia, New Zealand or India they would be awarded a Bachelor of Veterinary Science, Surgery or Medicine (BVS, BVSc, BVetMed or BVMS), and in Ireland graduates receive a Medicinae Veterinariae Baccalaureus (MVB).
In continental Europe, the degree of Doctor Medicinae Veterinariae (DMV, DrMedVet, Dr. med. vet., MVDr.) or Doctor Veterinariae Medicinae (DVM, DrVetMed, Dr. vet. med.) is granted.
The award of a bachelor's degree was previously commonplace in the United States, but the degree name and academic standards were upgraded to match the 'doctor' title used by graduates.
Comparatively few universities have veterinary schools that offer degrees which are accredited to qualify the graduates as registered vets. For example, there are 30 in the United States, 5 in Canada, 1 in New Zealand, 7 in Australia (4 of which offer degrees accredited by the American Veterinary Medical Association (AVMA)), and 8 in the United Kingdom (4 of which offer degrees accredited by the American Veterinary Medical Association (AVMA)).
Due to this scarcity of places for veterinary degrees, admission to veterinary school is competitive and requires extensive preparation. In the United States in 2007, approximately 5,750 applicants competed for the 2,650 seats in the 28 accredited veterinary schools, with an acceptance rate of 46%.
With competitive admission, many schools may place heavy emphasis and consideration on a candidate's veterinary and animal experience. Formal experience is a particular advantage to the applicant, often consisting of work with veterinarians or scientists in clinics, agribusiness, research, or some area of health science. Less formal experience is also helpful for the applicant to have, and this includes working with animals on a farm or ranch or at a stable or animal shelter and basic overall animal exposure.
In the United States, approximately 80% of admitted students are female. In the early history of veterinary medicine of the United States, most veterinarians were males. However, in the 1990s this ratio reached parity, and now it has been reversed.
Preveterinary courses should emphasize the sciences. Most veterinary schools typically require applicants to have taken one year equivalent classes in organic, inorganic chemistry, physics, general biology; and one semester of vertebrate embryology and biochemistry. Usually, the minimal mathematics requirement is college level calculus. Individual schools might require introduction to animal science, livestock judging, animal nutrition, cell biology, and genetics. However, due to the limited availability of these courses, many schools have removed these requirements to widen the pool of possible applicants.
Registration and licensing
Following academic education, most countries require a vet to be registered with the relevant governing body, and to maintain this license to practice.
According to the Bureau of Labor Statistics, veterinarians must be licensed to practice in the United States. Licensing entails passing an accredited program, a national exam, and a state exam. For instance, in the United States, a prospective vet must receive a passing grade on a national board examination, the North America Veterinary Licensing Exam. This exam must be completed over the course of eight hours, and consists of 360 multiple-choice questions, covering all aspects of veterinary medicine, as well as visual material designed to test diagnostic skills.
Postgraduate study
The percentage electing to undertake further study following registration in the United States has increased from 36.8% to 39.9% in 2008. About 25% of those or about 9% of graduates were accepted into traditional academic internships. Approximately 9% of veterinarians eventually board certify in one of 40 distinct specialties from 22 specialty organizations recognized by the AVMA American Board of Veterinary Specialties (ABVS).
ABVS recognized veterinary specialties
Source:
Curriculum comparison with human medicine
The first two-year curriculum in both veterinary and human medical schools are very similar in course names, but in certain subjects are relatively different in content. Considering the courses, the first two-year curriculum usually includes biochemistry, physiology, histology, anatomy, pharmacology, microbiology, epidemiology, pathology and hematology.
Some veterinary schools use the same biochemistry, histology, and microbiology books as human medical students; however, the course content is greatly supplemented to include the varied animal diseases and species differences. In the past, many veterinarians were trained in pharmacology using the same text books used by physicians. As the specialty of veterinary pharmacology has developed, more schools are using pharmacology textbooks written specifically for veterinarians. Veterinary physiology, anatomy, and histology is complex, as physiology often varies among species. Microbiology and virology of animals share the same foundation as human microbiology, but with grossly different disease manifestation and presentations. Epidemiology is focused on herd health and prevention of herd borne diseases and foreign animal diseases. Pathology, like microbiology and histology, is very diverse and encompasses many species and organ systems. Most veterinary schools have courses in small animal and large animal nutrition, often taken as electives in the clinical years or as part of the core curriculum in the first two years.
The final two-year curriculum is similar to that of human medicine only in clinical emphasis. A veterinary student must be well prepared to be a fully functional veterinarian on the day of graduation, competent in both surgery and medicine. The graduating veterinarian must be able to pass medical board examination and be prepared to enter clinical practice on the day of graduation, while most human medical doctors in the United States complete 3 to 5 years of post-doctoral residency before practicing medicine independently, usually in a very narrow and focused specialty. Many veterinarians do also complete a post-doctoral residency, but it is not nearly as common as it is in human medicine.
In the last years, curricula in both human and veterinary medicine have been adapted with the aim of incorporating competency-based teaching. Furthermore, the importance of institutionalized systematic teacher feedback has been recognized and tools such as clinical encounter cards are being implemented in clinical veterinary education.
Impact on human medicine
Some veterinarians pursue post-graduate training and enter research careers and have contributed to advances in many human and veterinary medical fields, including pharmacology and epidemiology. Research veterinarians were the first to isolate oncoviruses, Salmonella species, Brucella species, and various other pathogenic agents. Veterinarians were in the forefront in the effort to suppress malaria and yellow fever in the United States. Veterinarians identified the botulism disease-causing agent, developed propofol; a widely used anesthetic induction drug, produced an anticoagulant used to treat human heart disease, and developed surgical techniques for humans, such as hip-joint replacement, limb and organ transplants.
Occupational hazards
Veterinarians work with a wide variety of animal species typically in hospitals, clinics, labs, farms, and zoos. Veterinarians face many occupational hazards including zoonotic diseases, bites and scratches, hazardous drugs, needlestick injuries, ionizing radiation, and noise. According to the U.S. Department of Labor, 12% of workers in the veterinary services profession reported a work-related injury or illness in 2016.
Veterinary practices need a health and safety plan that addresses infection prevention and other hazards. Workplaces should utilize engineering controls, administrative controls, and personal protective equipment to keep their employees safe. PPE such as gloves, safety goggles, lab coats, and hearing protection should be readily available with mandatory training on proper usage. Raising awareness is the most important step in promoting workplace health and safety.
Biological and chemical hazards
Needlestick injuries are the most common accidents among veterinarians, but they are likely underreported. Needlesticks can result in hazardous drug or bloodborne-pathogen exposures.
Unlike human medical professionals, veterinarians receive minimal training on safe handling of hazardous drugs in school. Also, a large percentage of veterinarians are women of reproductive age and drug exposures put them at risk of infertility or other adverse health outcomes. Additionally, some antibiotics, steroids, and chemotherapy drugs are known to have negative effects on male fertility. The U.S. National Institute for Occupational Safety and Health has issued guidance on the safe handling of hazardous drugs for veterinary workers. Animal bites and scratches are another common injury in veterinary practice.
The close interactions with animals put veterinarians at increased risk of contracting zoonoses. A systematic review of veterinary students found that between 17% and 64% had acquired a zoonotic disease during their studies. The animal species, work setting, health and safety practices, and training can all affect the risk of injury and illness.
Physical hazards
Noise can be a prominent exposure, in which case a hearing loss prevention program may be recommended. A NIOSH study on kennel noise found that noise levels often exceeded OSHA's permissible exposure limit. Reducing noise is beneficial for animal and human health.
Psychosocial hazards
Veterinarians have high suicide rates in comparison to the general population. A study by the U.S. Centers for Disease Control and Prevention found that male veterinarians are 2.1 times and female veterinarians are 3.5 times as likely as the general population to die by suicide. Some reasons for this could be long hours, work overload, client expectations and complaints, poor remuneration, euthanasia procedures, and poor work-life balance. A survey of more than 11,000 vets found 9% had serious psychological distress, 31% experienced depressive episodes, and 17% had suicidal ideation. Online support groups, such as Not One More Vet, have been established to help veterinarians who may be experiencing suicidal thoughts. NOMV educates veterinarians and vet techs about other ways to help themselves with mental health. Another driver of stress can be student loan debt. A 2013 national survey found that average debt for veterinary medicine graduates was as high as $162,113. Veterinarian lifelong earning potential is less than a physician, so it can take a lot longer to break even.
In popular culture
Reality televisions shows featuring veterinarians include:
Bondi Vet, an Australian factual television series documenting the work of veterinary surgeon Chris Brown at the Bondi Junction Veterinary Hospital.
Dr. Oakley, Yukon Vet, about a Canadian veterinarian in the Yukon two of whose daughters assist her.
The Incredible Dr. Pol, a US veterinarian reality show. Produced by National Geographic Wild, a Disney channel. It follows the life of Dr. Jan Pol and Pol Veterinarian Service in Michigan.
E-Vet Interns (1998–2002), a US show filmed at Alameda East Veterinary Hospital in Denver, Colorado.
Emergency Vets, filmed at Alameda East Veterinary Hospital in Denver, Colorado.
Rookie Vets (2005), featuring students at Massey University in New Zealand.
Vet School Confidential (2001), following students at Michigan State University College of Veterinary Medicine in the US.
Vets in Practice (1997–2002), a British series.
Fictional works featuring a veterinarian as the main protagonist include:
James Herriot's series of books containing fictionalized stories of his career as a farm animal veterinarian in England, which was adapted as the BBC television series All Creatures Great and Small.
The Three Lives of Thomasina about Andrew MacDhui, a veterinarian in a village in Scotland.
The Doctor Dolittle series of children's books, which have thrice been adapted into movies, Doctor Dolittle (1967), Dr. Dolittle (1998), and Dolittle (2020).
The movie Beethoven, featuring the evil veterinarian Dr. Herman Varnick.
Veterinary malpractice
Most states in the US allow for malpractice lawsuit in case of death or injury to an animal from professional negligence. Usually the penalty is not greater than the value of the animal. Some states allow for punitive penalty, loss of companionship, and suffering, likely increasing the cost of veterinary malpractice insurance and the cost of veterinary care. Most veterinarians carry business, worker's compensation, and facility insurance to protect their clients and workers from injury inflicted by animals.
| Biology and health sciences | Health professionals | Health |
229167 | https://en.wikipedia.org/wiki/Bucephalus | Bucephalus | Bucephalus (; ; – June 326 BC) or Bucephalas, was the horse of Alexander the Great, and one of the most famous horses of classical antiquity. According to the Alexander Romance (1.15), the name "Bucephalus" literally means "ox-headed" (from and ), and supposedly comes from a brand (or scar) on the thigh of the horse that looked like an ox's head.
Ancient historical accounts state that Bucephalus's breed was that of the "best Thessalian strain", and that he died in what is now Punjab, Pakistan, after the Battle of the Hydaspes in 326 BC. Alexander was so grieved at the loss of his horse that he named one of the many cities he founded after him, as Alexandria Bucephalus.
Taming of Bucephalus
A massive creature with a massive head, Bucephalus is described as having a black coat with a large white star on his brow. He is also supposed to have had a "wall eye" (blue eye), and his breeding was that of the "best Thessalian strain".
Plutarch says that in 344 BC, at twelve or thirteen years of age, Alexander of Macedonia won the horse by making a wager with his father: a horse dealer named Philonicus the Thessalian offered Bucephalus to King Philip II for the remarkably high sum of 13 talents. Because no one could tame the animal, Philip was not interested. However, Alexander was, and he offered to pay himself should he fail.
Alexander was given a chance and surprised all by subduing the horse. He spoke soothingly to the horse and turned its head toward the sun so that it could no longer see its own shadow, which had been the cause of its distress. Dropping his fluttering cloak as well, Alexander successfully tamed the horse. Plutarch says that the incident so impressed Philip that he told the boy, "O my son, look thee out a kingdom equal to and worthy of thyself, for Macedonia is too little for thee." Philip's speech strikes the only false note in the anecdote, according to A. R. Anderson, who noted his words as the embryo of the legend fully developed in the History of Alexander the Great I.15, 17.
The Alexander Romance presents a mythic variant of Bucephalus's origin. In this tale, the colt, whose heroic attributes surpassed even those of Pegasus, is bred and presented to Philip on his own estates. The mythic attributes of the animal are further reinforced in the romance by the Delphic Oracle who tells Philip that the destined king of the world will be the one who rides Bucephalus, a horse with the mark of the ox's head on his haunch.
Alexander and Bucephalus
As one of his chargers, Bucephalus served Alexander in numerous battles.
The value which Alexander placed on Bucephalus emulated his hero and supposed ancestor Achilles, who claimed that his horses were "known to excel all others—for they are immortal. Poseidon gave them to my father Peleus, who in his turn gave them to me."
Arrian states, with Onesicritus as his source, that Bucephalus died at the age of thirty. Other sources, however, give as the cause of death not old age or weariness, but fatal injuries at the Battle of the Hydaspes (June 326 BC), in which Alexander's army defeated King Porus. Alexander promptly founded a city, Bucephala, in honour of his horse. It was on the west bank of the Hydaspes river (modern-day Jhelum in Pakistan). The modern-day town of Jalalpur Sharif, outside Jhelum, is said to be where Bucephalus is buried.
The legend of Bucephalus grew in association with that of Alexander, beginning with the fiction that they were born simultaneously: some of the later versions of the Alexander Romance also synchronized the hour of their death. The Bucephalus appears in almost all versions of the Armenian Alexander Romance, and visual illustrations in the surviving manuscripts of this text sometimes represent scenes with the Bucephalus.
| Biology and health sciences | Individual animals | Animals |
229173 | https://en.wikipedia.org/wiki/Plecoptera | Plecoptera | Plecoptera is an order of insects, commonly known as stoneflies. Some 3,500 species are described worldwide, with new species still being discovered. Stoneflies are found worldwide, except Antarctica. Stoneflies are believed to be one of the most primitive groups of Neoptera, with close relatives identified from the Carboniferous and Lower Permian geological periods, while true stoneflies are known from fossils only a bit younger. Their modern diversity, however, apparently is of Mesozoic origin.
Plecoptera are found in both the Southern and Northern Hemispheres, and the populations are quite distinct, although the evolutionary evidence suggests species may have crossed the equator on a number of occasions before once again becoming geographically isolated.
All species of Plecoptera are intolerant of water pollution, and the presence of their nymphs in a stream or still water is usually an indicator of good or excellent water quality.
Description and ecology
Stoneflies have a generalized anatomy, with few specialized features compared to other insects. They have simple mouthparts with chewing mandibles, long, multiple-segmented antennae, large compound eyes, and two or three ocelli. The legs are robust, with each ending in two claws. The abdomen is relatively soft, and may include remnants of the nymphal gills even in the adult. Both nymphs and adults have long, paired cerci projecting from the tip of their abdomens.
The name "Plecoptera" literally means "braided-wings", from the Ancient Greek plekein (πλέκειν, "to braid") and pteryx (πτέρυξ, "wing"). This refers to the complex venation of their two pairs of wings, which are membranous and fold flat over their backs. Stoneflies are generally not strong fliers, and some species are entirely wingless.
A few wingless species, such as the Lake Tahoe benthic stonefly ("Capnia" lacustra) or Baikaloperla, are the only known insects, perhaps with the exception of Halobates, that are exclusively aquatic from birth to death. Some true water bugs (Nepomorpha) may also be fully aquatic for their entire lives, but can leave the water to travel.
The nymphs (technically, "naiads") are aquatic and live in the benthic zone of well-oxygenated lakes and streams. A few species found in New Zealand and nearby islands have terrestrial nymphs, but even these inhabit only very moist environments. The nymphs physically resemble wingless adults, but often have external gills, which may be present on almost any part of the body. Nymphs can acquire oxygen via diffusing through the exoskeleton, or through gills located on behind the head, on the thorax, or around the anus. Due to their nymph's requirement for well oxygenated water, the species is very sensitive to water pollution. This makes them important indicators for water quality. Most species are herbivorous as nymphs, feeding on submerged leaves and benthic algae, but many are hunters of other aquatic arthropods.
Life cycle
The female can lay up to one thousand eggs. It will fly over the water and drop the eggs in the water. It also may hang on a rock or branch. Eggs are covered in a sticky coating which allows them to adhere to rocks without being swept away by swift currents. The eggs typically take two to three weeks to hatch, but some species undergo diapause, with the eggs remaining dormant throughout a dry season, and hatching only when conditions are suitable.
The insects remain in the nymphal form for one to four years, depending on species, and undergo from 12 to 36 molts before emerging and becoming terrestrial as adults. Before becoming adults, nymphs will leave the water, attach to a fixed surface and molt one last time.
The adults generally only survive for a few weeks, and emerge only during specific times of the year when resources are optimal. Some do not feed at all, but those that do are herbivorous. Adults are not strong fliers and generally stay near the stream or lake they hatched from.
Phylogeny
A summary of the phylogeny of stoneflies is shown below. While the Antarctoperlaria, Arctoperlaria, Euholognatha, Systellognatha are well supported, several further relationships are disputed. Some families have only been analyzed by one study (indicated by dashed lines) and their placement may change in the near future.
| Biology and health sciences | Insects: General | Animals |
229253 | https://en.wikipedia.org/wiki/Dieffenbachia | Dieffenbachia | Dieffenbachia , commonly known as dumb cane or leopard lily, is a genus of tropical flowering plants in the family Araceae. It is native to the New World Tropics from Mexico and the West Indies south to Argentina. Some species are widely cultivated as ornamental plants, especially as houseplants, and have become naturalized on a few tropical islands.
Dieffenbachia is a perennial herbaceous plant with straight stem, simple and alternate leaves containing white spots and flecks, making it attractive as indoor foliage. Species in this genus are popular as houseplants because of their tolerance of shade. The English names, dumb cane and mother-in-law's tongue (also used for Sansevieria species) refer to the poisoning effect of raphides, which can cause temporary inability to speak. Dieffenbachia was named by Heinrich Wilhelm Schott, director of the Botanical Gardens in Vienna, to honor his head gardener Joseph Dieffenbach (1790–1863).
Species
The World Checklist of Selected Plant Families lists the following species:
Dieffenbachia aglaonematifolia Engl. – Brazil, Paraguay; Corrientes + Misiones Provinces of Argentina
Dieffenbachia antioquensis Linden ex Rafarin – Colombia
Dieffenbachia aurantiaca Engl – Costa Rica, Panama
Dieffenbachia beachiana Croat & Grayum – Costa Rica, Panama
Dieffenbachia bowmannii Carrière – Colombia, northwestern Brazil
Dieffenbachia brittonii Engl. – Colombia
Dieffenbachia burgeri Croat & Grayum – Costa Rica
Dieffenbachia cannifolia Engl. – Colombia, Ecuador, Peru
Dieffenbachia concinna Croat & Grayum – Costa Rica, Nicaragua
Dieffenbachia copensis Croat – Panama
Dieffenbachia cordata Engl. – Peru
Dieffenbachia costata Klotzsch ex Schott – Colombia, Ecuador, Peru
Dieffenbachia crebripistillata Croat – Panama
Dieffenbachia daguensis Engl. – Colombia, Ecuador
Dieffenbachia davidsei Croat & Grayum – Costa Rica
Dieffenbachia duidae (Steyerm.) G.S.Bunting – Venezuela, Guyana
Dieffenbachia elegans A.M.E.Jonker & Jonker – Bolivia, northwestern Brazil, the Guianas
Dieffenbachia enderi Engl. – Colombia
Dieffenbachia fortunensis Croat – Panama
Dieffenbachia fosteri Croat – Panama
Dieffenbachia fournieri N.E.Br. – Colombia
Dieffenbachia galdamesiae Croat – Panama
Dieffenbachia gracilis Huber – Peru, northwestern Brazil
Dieffenbachia grayumiana Croat – Costa Rica, Panama, Colombia
Dieffenbachia hammelii Croat & Grayum – Costa Rica, Nicaragua
Dieffenbachia herthae Diels – Ecuador
Dieffenbachia horichii Croat & Grayum – Costa Rica
Dieffenbachia humilis Poepp. – Bolivia, Peru, Ecuador, northwestern Brazil, the Guianas
Dieffenbachia imperialis Linden & André – Peru
Dieffenbachia isthmia Croat – Panama
Dieffenbachia killipii Croat – Panama
Dieffenbachia lancifolia Linden & André – Colombia
Dieffenbachia leopoldii W.Bull – Colombia
Dieffenbachia longispatha Engl. & K.Krause – Panama, Colombia
Dieffenbachia lutheri Croat – Panama
Dieffenbachia macrophylla Poepp. – Peru
Dieffenbachia meleagris L.Linden & Rodigas – Ecuador
Dieffenbachia nitidipetiolata Croat & Grayum – Panama
Dieffenbachia obliqua Poepp. – Peru
Dieffenbachia obscurinervia Croat – Panama
Dieffenbachia oerstedii Schott – southern Mexico (Veracruz, Tabasco, Campeche, Oaxaca, Chiapas), Central America (all 7 countries), Colombia
Dieffenbachia olbia L.Linden & Rodigas – Peru
Dieffenbachia paludicola N.E.Br. ex Gleason – northwestern Brazil, the Guianas, Venezuela
Dieffenbachia panamensis Croat – Panama
Dieffenbachia parlatorei Linden & André – Colombia, Venezuela
Dieffenbachia parvifolia Engl. – northwestern Brazil, Bolivia, Ecuador, Peru, Venezuela
Dieffenbachia pittieri Engl. & K.Krause – Panama
Dieffenbachia seguine (Jacq.) Schott – West Indies, south to Brazil and Bolivia (syn. Dieffenbachia maculata, Dieffenbachia picta)
Dieffenbachia shuttleworthiana Regel – Colombia
Dieffenbachia standleyi Croat – Honduras
Dieffenbachia tonduzii Croat & Grayum – Nicaragua, Costa Rica, Panama, Colombia, Ecuador
Dieffenbachia weberbaueri Engl. – Peru
Dieffenbachia weirii Berk. – Colombia
Dieffenbachia wendlandii Schott – southern Mexico (Querétaro, Veracruz, Oaxaca, Chiapas) south to Panama
Dieffenbachia williamsii Croat – Bolivia
Dieffenbachia wurdackii Croat – Peru
Ecology
In a survey that began in 1998, researchers in Costa Rica noticed that the strawberry poison frog Oophaga pumilio, deposited almost all (89%) of their tadpoles on the leaf axils of Dieffenbachia. As a result, the frog population fluctuated with the abundance of Dieffenbachia, especially in secondary forests
. A majority of the plants were eradicated by 2012 when the surveyors returned to the same area, with only 28% of 2002 plant numbers remaining. Researchers concluded that the reason for the rapid decline in Dieffenbachia was due to increased abundance of the collared peccary Dicotyles tajacu in the La Selva Biological Station research area; a small pig-like animal that feeds on Dieffenbachia and other plants.
Cultivation
With a minimum temperature of , dieffenbachia must be grown indoors in temperate areas. They need light, but filtered sunlight through a window is usually sufficient. They also need moderately moist soil, which should be regularly fertilized with an appropriate houseplant fertilizer. Leaves will periodically roll up and fall off to make way for new leaves. Yellowing of the leaves is generally a sign of problematic conditions, such as a nutrient deficiency in the soil. Dieffenbachia respond well to hot temperatures and dry climates.
Dieffenbachia prefer medium sunlight, moderately dry soil and average home temperatures of . Most require water about twice a week.
As Dieffenbachia seguine comes from the tropical rain forest, it prefers to have moisture at its roots, as it grows all the time, it needs constant water, but with loose well aerated soils.
The cultivars 'Camille' and 'Tropic Snow' have gained the Royal Horticultural Society's Award of Garden Merit.
Toxicity
The cells of the Dieffenbachia plant contain needle-shaped calcium oxalate crystals called raphides. If a leaf is chewed, these crystals can cause a temporary burning sensation and erythema. In rare cases, edema of tissues exposed to the plant has been reported. Mastication and ingestion generally result in only mild symptoms. With both children and pets, contact with Dieffenbachia (typically from chewing) can cause a host of unpleasant symptoms, including intense numbing, oral irritation, excessive drooling, and localized swelling. However, these effects are rarely life-threatening. In most cases, symptoms are mild, and can be successfully treated with analgesic agents, antihistamines, or medical charcoal.
Severe cases can occur if Dieffenbachia makes prolonged contact with oral mucosal tissue. In such cases, symptoms generally include severe pain which can last for several days to weeks. Hospitalization may be necessary if prolonged contact is made with the throat, in which severe swelling has the potential to affect breathing.
Gastric evacuation or lavage is "seldom" indicated. In patients with exposure to toxic plants, 70% are children younger than 5 years.
Stories that Dieffenbachia is a deadly poison are urban legends.
| Biology and health sciences | Alismatales | Plants |
229296 | https://en.wikipedia.org/wiki/Pons | Pons | The pons (from Latin , "bridge") is part of the brainstem that in humans and other mammals, lies inferior to the midbrain, superior to the medulla oblongata and anterior to the cerebellum.
The pons is also called the pons Varolii ("bridge of Varolius"), after the Italian anatomist and surgeon Costanzo Varolio (1543–75). This region of the brainstem includes neural pathways and tracts that conduct signals from the brain down to the cerebellum and medulla, and tracts that carry the sensory signals up into the thalamus.
Structure
The pons in humans measures about in length. It is the part of the brainstem situated between the midbrain and the medulla oblongata. The horizontal medullopontine sulcus demarcates the boundary between the pons and medulla oblongata on the ventral aspect of the brainstem, and the roots of cranial nerves VI/VII/VIII emerge from the brainstem along this groove. The junction of pons, medulla oblongata, and cerebellum forms the cerebellopontine angle. The superior pontine sulcus separates the pons from the midbrain. Posteriorly, the pons curves on either side into a middle cerebellar peduncle.
A cross-section of the pons divides it into a ventral and a dorsal area. The ventral pons is known as the basilar part, and the dorsal pons is known as the pontine tegmentum.
The ventral aspect of the pons faces the clivus, with the pontine cistern intervening between the two structures. The ventral surface of the pons features a midline basilar sulcus along which the basilar artery may or may not course. There is a bulge to either side of the basilar sulcus, created by the pontine nuclei that are interweaved amid the descending fibres within the substance of the pons. The superior cerebellar artery winds around the upper margin of the pons.
Vasculature
Most of the pons is supplied by the pontine arteries, which arise from the basilar artery. A smaller portion of the pons is supplied by the anterior and posterior inferior cerebellar arteries.
Development
During embryonic development, the metencephalon develops from the rhombencephalon and gives rise to two structures: the pons and the cerebellum. The alar plate produces sensory neuroblasts, which will give rise to the solitary nucleus and its special visceral afferent (SVA) column; the cochlear and vestibular nuclei, which form the special somatic afferent (SSA) fibers of the vestibulocochlear nerve, the spinal and principal trigeminal nerve nuclei, which form the general somatic afferent column (GSA) of the trigeminal nerve, and the pontine nuclei which relays to the cerebellum.
Basal plate neuroblasts give rise to the abducens nucleus, which forms the general somatic efferent fibers (GSE); the facial and motor trigeminal nuclei, which form the special visceral efferent (SVE) column, and the superior salivatory nucleus, which forms the general visceral efferent fibers (GVE) of the facial nerve.
Nuclei
A number of cranial nerve nuclei are present in the pons:
mid-pons: the principal sensory nucleus of the trigeminal nerve (V)
mid-pons: the motor nucleus for the trigeminal nerve (V)
lower down in the pons: abducens nucleus (VI)
lower down in the pons: facial nerve nucleus (VII)
lower down in the pons: vestibulocochlear nuclei (vestibular nuclei and cochlear nuclei) (VIII)
Function
Functions of these four cranial nerves (V-VIII) include regulation of respiration, control of involuntary actions, sensory roles in hearing, equilibrium, and taste, and in facial sensations such as touch and pain, as well as motor roles in eye movement, facial expressions, chewing, swallowing, and the secretion of saliva and tears.
The pons contains nuclei that relay signals from the forebrain to the cerebellum, along with nuclei that deal primarily with sleep, respiration, swallowing, bladder control, hearing, equilibrium, taste, eye movement, facial expressions, facial sensation, and posture.
Within the pons is the pneumotaxic center consisting of the subparabrachial and the medial parabrachial nuclei. This center regulates the change from inhalation to exhalation.
The pons is implicated in sleep paralysis, and may also play a role in generating dreams.
Clinical significance
Central pontine myelinolysis is a demyelinating disease that causes difficulty with sense of balance, walking, sense of touch, swallowing and speaking. In a clinical setting, it is often associated with transplant or rapid correction of blood sodium. Undiagnosed, it can lead to death or locked-in syndrome.
Other animals
Evolution
The pons first evolved as an offshoot of the medullary reticular formation. Since lampreys possess a pons, it has been argued that it must have evolved as a region distinct from the medulla by the time the first agnathans appeared, 525 million years ago.
Additional images
| Biology and health sciences | Nervous system | Biology |
229470 | https://en.wikipedia.org/wiki/Haikouichthys | Haikouichthys | Haikouichthys is an extinct genus of craniate (animals with notochords and distinct heads) that lived 518 million years ago, during the Cambrian explosion of multicellular life. The type species, Haikouichthys ercaicunensis, was first described in 1999. Haikouichthys had a defined skull and other characteristics that have led paleontologists to label it a true craniate, and even to be popularly characterized as one of the earliest fishes. More than 500 specimens were referred to this taxon and phylogenetic analyses indicates that the animal is probably a basal stem-craniate. Some researchers have considered Haikouichthys to be synonymous with the other primitive chordate Myllokunmingia, but subsequent studies led by the British paleontologist Simon Conway Morris identified both genera to be distinct, separate taxa on the basis of different gill arrangement, the absence of branchial rays in Myllokunmingia and the myomeres having a more acute shape in Haikouichthys.
Description
Haikouichthys is about long and is narrower than Myllokunmingia, another putative chordate that comes from the same beds. The holotype of Haikouichthys ercaicunensis was found in the Yuanshan member of the Qiongzhusi Formation in the 'Eoredlichia' Zone near Ercai Village in the Haikou Subdistrict (not to be confused with the city of Haikou in Hainan) of Xishan, Kunming, hence its name, which means "Haikou fish from Ercaicun". The fossil was recovered among the Chengjiang fauna, in one of a series of Lagerstätten sites where thousands of exquisitely preserved soft-bodied fossils have already been found. Following the discovery of the holotype, additional Lower Cambrian fossils of Haikouichthys ercaicunensis have been discovered.
Researchers have identified eyes, cranial cartilages, at least six to nine gill arches, possible nasal sacs and otic capsules from its head. It is likely that the brain of Haikouichthys had the same major brain divisions found in extant vertebrates. The describers of this taxon initially reported its potential notochord, though some researchers consider this claim to be uncertain. Still, numerous segments (myomeres) with rear directed chevrons in the tail indicate that Haikouichthys was indeed a chordate, and complete dorsal, ventral and caudal fins were also found in its specimens. The fin radials of Haikouichthys show similarity to those of hagfish and lampreys, and they seem to angle "forward" toward the end thought on the basis of internal structures to be the head. There are 13 circular structures along the bottom that may be gonads, slime organs, or something else entirely.
| Biology and health sciences | Prehistoric agnathae and early chordates | Animals |
229553 | https://en.wikipedia.org/wiki/Hooke%27s%20law | Hooke's law | In physics, Hooke's law is an empirical law which states that the force () needed to extend or compress a spring by some distance () scales linearly with respect to that distance—that is, where is a constant factor characteristic of the spring (i.e., its stiffness), and is small compared to the total possible deformation of the spring. The law is named after 17th-century British physicist Robert Hooke. He first stated the law in 1676 as a Latin anagram. He published the solution of his anagram in 1678 as: ("as the extension, so the force" or "the extension is proportional to the force"). Hooke states in the 1678 work that he was aware of the law since 1660.
Hooke's equation holds (to some extent) in many other situations where an elastic body is deformed, such as wind blowing on a tall building, and a musician plucking a string of a guitar. An elastic body or material for which this equation can be assumed is said to be linear-elastic or Hookean.
Hooke's law is only a first-order linear approximation to the real response of springs and other elastic bodies to applied forces. It must eventually fail once the forces exceed some limit, since no material can be compressed beyond a certain minimum size, or stretched beyond a maximum size, without some permanent deformation or change of state. Many materials will noticeably deviate from Hooke's law well before those elastic limits are reached.
On the other hand, Hooke's law is an accurate approximation for most solid bodies, as long as the forces and deformations are small enough. For this reason, Hooke's law is extensively used in all branches of science and engineering, and is the foundation of many disciplines such as seismology, molecular mechanics and acoustics. It is also the fundamental principle behind the spring scale, the manometer, the galvanometer, and the balance wheel of the mechanical clock.
The modern theory of elasticity generalizes Hooke's law to say that the strain (deformation) of an elastic object or material is proportional to the stress applied to it. However, since general stresses and strains may have multiple independent components, the "proportionality factor" may no longer be just a single real number, but rather a linear map (a tensor) that can be represented by a matrix of real numbers.
In this general form, Hooke's law makes it possible to deduce the relation between strain and stress for complex objects in terms of intrinsic properties of the materials they are made of. For example, one can deduce that a homogeneous rod with uniform cross section will behave like a simple spring when stretched, with a stiffness directly proportional to its cross-section area and inversely proportional to its length.
Formal definition
Linear springs
Consider a simple helical spring that has one end attached to some fixed object, while the free end is being pulled by a force whose magnitude is . Suppose that the spring has reached a state of equilibrium, where its length is not changing anymore. Let be the amount by which the free end of the spring was displaced from its "relaxed" position (when it is not being stretched). Hooke's law states that or, equivalently,
where is a positive real number, characteristic of the spring. A spring with spaces between the coils can be compressed, and the same formula holds for compression, with and both negative in that case.
According to this formula, the graph of the applied force as a function of the displacement will be a straight line passing through the origin, whose slope is .
Hooke's law for a spring is also stated under the convention that is the restoring force exerted by the spring on whatever is pulling its free end. In that case, the equation becomes since the direction of the restoring force is opposite to that of the displacement.
Torsional springs
The torsional analog of Hooke's law applies to torsional springs. It states that the torque (τ) required to rotate an object is directly proportional to the angular displacement (θ) from the equilibrium position. It describes the relationship between the torque applied to an object and the resulting angular deformation due to torsion. Mathematically, it can be expressed as:
Where:
τ is the torque measured in Newton-meters or N·m.
k is the torsional constant (measured in N·m/radian), which characterizes the stiffness of the torsional spring or the resistance to angular displacement.
θ is the angular displacement (measured in radians) from the equilibrium position.
Just as in the linear case, this law shows that the torque is proportional to the angular displacement, and the negative sign indicates that the torque acts in a direction opposite to the angular displacement, providing a restoring force to bring the system back to equilibrium.
General "scalar" springs
Hooke's spring law usually applies to any elastic object, of arbitrary complexity, as long as both the deformation and the stress can be expressed by a single number that can be both positive and negative.
For example, when a block of rubber attached to two parallel plates is deformed by shearing, rather than stretching or compression, the shearing force and the sideways displacement of the plates obey Hooke's law (for small enough deformations).
Hooke's law also applies when a straight steel bar or concrete beam (like the one used in buildings), supported at both ends, is bent by a weight placed at some intermediate point. The displacement in this case is the deviation of the beam, measured in the transversal direction, relative to its unloaded shape.
Vector formulation
In the case of a helical spring that is stretched or compressed along its axis, the applied (or restoring) force and the resulting elongation or compression have the same direction (which is the direction of said axis). Therefore, if and are defined as vectors, Hooke's equation still holds and says that the force vector is the elongation vector multiplied by a fixed scalar.
General tensor form
Some elastic bodies will deform in one direction when subjected to a force with a different direction. One example is a horizontal wood beam with non-square rectangular cross section that is bent by a transverse load that is neither vertical nor horizontal. In such cases, the magnitude of the displacement will be proportional to the magnitude of the force , as long as the direction of the latter remains the same (and its value is not too large); so the scalar version of Hooke's law will hold. However, the force and displacement vectors will not be scalar multiples of each other, since they have different directions. Moreover, the ratio between their magnitudes will depend on the direction of the vector .
Yet, in such cases there is often a fixed linear relation between the force and deformation vectors, as long as they are small enough. Namely, there is a function from vectors to vectors, such that , and for any real numbers , and any displacement vectors , . Such a function is called a (second-order) tensor.
With respect to an arbitrary Cartesian coordinate system, the force and displacement vectors can be represented by 3 × 1 matrices of real numbers. Then the tensor connecting them can be represented by a 3 × 3 matrix of real coefficients, that, when multiplied by the displacement vector, gives the force vector:
That is, for . Therefore, Hooke's law can be said to hold also when and are vectors with variable directions, except that the stiffness of the object is a tensor , rather than a single real number .
Hooke's law for continuous media
The stresses and strains of the material inside a continuous elastic material (such as a block of rubber, the wall of a boiler, or a steel bar) are connected by a linear relationship that is mathematically similar to Hooke's spring law, and is often referred to by that name.
However, the strain state in a solid medium around some point cannot be described by a single vector. The same parcel of material, no matter how small, can be compressed, stretched, and sheared at the same time, along different directions. Likewise, the stresses in that parcel can be at once pushing, pulling, and shearing.
In order to capture this complexity, the relevant state of the medium around a point must be represented by two-second-order tensors, the strain tensor (in lieu of the displacement ) and the stress tensor (replacing the restoring force ). The analogue of Hooke's spring law for continuous media is then where is a fourth-order tensor (that is, a linear map between second-order tensors) usually called the stiffness tensor or elasticity tensor. One may also write it as where the tensor , called the compliance tensor, represents the inverse of said linear map.
In a Cartesian coordinate system, the stress and strain tensors can be represented by 3 × 3 matrices
Being a linear mapping between the nine numbers and the nine numbers , the stiffness tensor is represented by a matrix of real numbers . Hooke's law then says that
where .
All three tensors generally vary from point to point inside the medium, and may vary with time as well. The strain tensor merely specifies the displacement of the medium particles in the neighborhood of the point, while the stress tensor specifies the forces that neighboring parcels of the medium are exerting on each other. Therefore, they are independent of the composition and physical state of the material. The stiffness tensor , on the other hand, is a property of the material, and often depends on physical state variables such as temperature, pressure, and microstructure.
Due to the inherent symmetries of , , and , only 21 elastic coefficients of the latter are independent. This number can be further reduced by the symmetry of the material: 9 for an orthorhombic crystal, 5 for an hexagonal structure, and 3 for a cubic symmetry. For isotropic media (which have the same physical properties in any direction), can be reduced to only two independent numbers, the bulk modulus and the shear modulus , that quantify the material's resistance to changes in volume and to shearing deformations, respectively.
Analogous laws
Since Hooke's law is a simple proportionality between two quantities, its formulas and consequences are mathematically similar to those of many other physical laws, such as those describing the motion of fluids, or the polarization of a dielectric by an electric field.
In particular, the tensor equation relating elastic stresses to strains is entirely similar to the equation relating the viscous stress tensor and the strain rate tensor in flows of viscous fluids; although the former pertains to static stresses (related to amount of deformation) while the latter pertains to dynamical stresses (related to the rate of deformation).
Units of measurement
In SI units, displacements are measured in meters (m), and forces in newtons (N or kg·m/s2). Therefore, the spring constant , and each element of the tensor , is measured in newtons per meter (N/m), or kilograms per second squared (kg/s2).
For continuous media, each element of the stress tensor is a force divided by an area; it is therefore measured in units of pressure, namely pascals (Pa, or N/m2, or kg/(m·s2). The elements of the strain tensor are dimensionless (displacements divided by distances). Therefore, the entries of are also expressed in units of pressure.
General application to elastic materials
Objects that quickly regain their original shape after being deformed by a force, with the molecules or atoms of their material returning to the initial state of stable equilibrium, often obey Hooke's law.
Hooke's law only holds for some materials under certain loading conditions. Steel exhibits linear-elastic behavior in most engineering applications; Hooke's law is valid for it throughout its elastic range (i.e., for stresses below the yield strength). For some other materials, such as aluminium, Hooke's law is only valid for a portion of the elastic range. For these materials a proportional limit stress is defined, below which the errors associated with the linear approximation are negligible.
Rubber is generally regarded as a "non-Hookean" material because its elasticity is stress dependent and sensitive to temperature and loading rate.
Generalizations of Hooke's law for the case of large deformations is provided by models of neo-Hookean solids and Mooney–Rivlin solids.
Derived formulae
Tensional stress of a uniform bar
A rod of any elastic material may be viewed as a linear spring. The rod has length and cross-sectional area . Its tensile stress is linearly proportional to its fractional extension or strain by the modulus of elasticity :
The modulus of elasticity may often be considered constant. In turn,
(that is, the fractional change in length), and since
it follows that:
The change in length may be expressed as
Spring energy
The potential energy stored in a spring is given by which comes from adding up the energy it takes to incrementally compress the spring. That is, the integral of force over displacement. Since the external force has the same general direction as the displacement, the potential energy of a spring is always non-negative. Substituting gives
This potential can be visualized as a parabola on the -plane such that . As the spring is stretched in the positive -direction, the potential energy increases parabolically (the same thing happens as the spring is compressed). Since the change in potential energy changes at a constant rate:
Note that the change in the change in is constant even when the displacement and acceleration are zero.
Relaxed force constants (generalized compliance constants)
Relaxed force constants (the inverse of generalized compliance constants) are uniquely defined for molecular systems, in contradistinction to the usual "rigid" force constants, and thus their use allows meaningful correlations to be made between force fields calculated for reactants, transition states, and products of a chemical reaction. Just as the potential energy can be written as a quadratic form in the internal coordinates, so it can also be written in terms of generalized forces. The resulting coefficients are termed compliance constants. A direct method exists for calculating the compliance constant for any internal coordinate of a molecule, without the need to do the normal mode analysis. The suitability of relaxed force constants (inverse compliance constants) as covalent bond strength descriptors was demonstrated as early as 1980. Recently, the suitability as non-covalent bond strength descriptors was demonstrated too.
Harmonic oscillator
A mass attached to the end of a spring is a classic example of a harmonic oscillator. By pulling slightly on the mass and then releasing it, the system will be set in sinusoidal oscillating motion about the equilibrium position. To the extent that the spring obeys Hooke's law, and that one can neglect friction and the mass of the spring, the amplitude of the oscillation will remain constant; and its frequency will be independent of its amplitude, determined only by the mass and the stiffness of the spring:
This phenomenon made possible the construction of accurate mechanical clocks and watches that could be carried on ships and people's pockets.
Rotation in gravity-free space
If the mass were attached to a spring with force constant and rotating in free space, the spring tension () would supply the required centripetal force ():
Since and , then:
Given that , this leads to the same frequency equation as above:
Linear elasticity theory for continuous media
Isotropic materials
Isotropic materials are characterized by properties which are independent of direction in space. Physical equations involving isotropic materials must therefore be independent of the coordinate system chosen to represent them. The strain tensor is a symmetric tensor. Since the trace of any tensor is independent of any coordinate system, the most complete coordinate-free decomposition of a symmetric tensor is to represent it as the sum of a constant tensor and a traceless symmetric tensor. Thus in index notation:
where is the Kronecker delta. In direct tensor notation:
where is the second-order identity tensor.
The first term on the right is the constant tensor, also known as the volumetric strain tensor, and the second term is the traceless symmetric tensor, also known as the deviatoric strain tensor or shear tensor.
The most general form of Hooke's law for isotropic materials may now be written as a linear combination of these two tensors:
where is the bulk modulus and is the shear modulus.
Using the relationships between the elastic moduli, these equations may also be expressed in various other ways. A common form of Hooke's law for isotropic materials, expressed in direct tensor notation, is
where and are the Lamé constants, is the second-rank identity tensor, and I is the symmetric part of the fourth-rank identity tensor. In index notation:
The inverse relationship is
Therefore, the compliance tensor in the relation is
In terms of Young's modulus and Poisson's ratio, Hooke's law for isotropic materials can then be expressed as
This is the form in which the strain is expressed in terms of the stress tensor in engineering. The expression in expanded form is
where is Young's modulus and is Poisson's ratio. (See 3-D elasticity).
In matrix form, Hooke's law for isotropic materials can be written as
where is the engineering shear strain. The inverse relation may be written as
which can be simplified thanks to the Lamé constants:
In vector notation this becomes
where is the identity tensor.
Plane stress
Under plane stress conditions, . In that case Hooke's law takes the form
In vector notation this becomes
The inverse relation is usually written in the reduced form
Plane strain
Under plane strain conditions, . In this case Hooke's law takes the form
Anisotropic materials
The symmetry of the Cauchy stress tensor () and the generalized Hooke's laws () implies that . Similarly, the symmetry of the infinitesimal strain tensor implies that . These symmetries are called the minor symmetries of the stiffness tensor c. This reduces the number of elastic constants from 81 to 36.
If in addition, since the displacement gradient and the Cauchy stress are work conjugate, the stress–strain relation can be derived from a strain energy density functional (), then
The arbitrariness of the order of differentiation implies that . These are called the major symmetries of the stiffness tensor. This reduces the number of elastic constants from 36 to 21. The major and minor symmetries indicate that the stiffness tensor has only 21 independent components.
Matrix representation (stiffness tensor)
It is often useful to express the anisotropic form of Hooke's law in matrix notation, also called Voigt notation. To do this we take advantage of the symmetry of the stress and strain tensors and express them as six-dimensional vectors in an orthonormal coordinate system () as
Then the stiffness tensor (c) can be expressed as
and Hooke's law is written as
Similarly the compliance tensor (s) can be written as
Change of coordinate system
If a linear elastic material is rotated from a reference configuration to another, then the material is symmetric with respect to the rotation if the components of the stiffness tensor in the rotated configuration are related to the components in the reference configuration by the relation
where are the components of an orthogonal rotation matrix . The same relation also holds for inversions.
In matrix notation, if the transformed basis (rotated or inverted) is related to the reference basis by
then
In addition, if the material is symmetric with respect to the transformation then
Orthotropic materials
Orthotropic materials have three orthogonal planes of symmetry. If the basis vectors () are normals to the planes of symmetry then the coordinate transformation relations imply that
The inverse of this relation is commonly written as
where
is the Young's modulus along axis
is the shear modulus in direction on the plane whose normal is in direction
is the Poisson's ratio that corresponds to a contraction in direction when an extension is applied in direction .
Under plane stress conditions, , Hooke's law for an orthotropic material takes the form
The inverse relation is
The transposed form of the above stiffness matrix is also often used.
Transversely isotropic materials
A transversely isotropic material is symmetric with respect to a rotation about an axis of symmetry. For such a material, if is the axis of symmetry, Hooke's law can be expressed as
More frequently, the axis is taken to be the axis of symmetry and the inverse Hooke's law is written as
Universal elastic anisotropy index
To grasp the degree of anisotropy of any class, a universal elastic anisotropy index (AU) was formulated. It replaces the Zener ratio, which is suited for cubic crystals.
Thermodynamic basis
Linear deformations of elastic materials can be approximated as adiabatic. Under these conditions and for quasistatic processes the first law of thermodynamics for a deformed body can be expressed as
where is the increase in internal energy and is the work done by external forces. The work can be split into two terms
where is the work done by surface forces while is the work done by body forces. If is a variation of the displacement field in the body, then the two external work terms can be expressed as
where is the surface traction vector, is the body force vector, represents the body and represents its surface. Using the relation between the Cauchy stress and the surface traction, (where is the unit outward normal to ), we have
Converting the surface integral into a volume integral via the divergence theorem gives
Using the symmetry of the Cauchy stress and the identity
we have the following
From the definition of strain and from the equations of equilibrium we have
Hence we can write
and therefore the variation in the internal energy density is given by
An elastic material is defined as one in which the total internal energy is equal to the potential energy of the internal forces (also called the elastic strain energy). Therefore, the internal energy density is a function of the strains, and the variation of the internal energy can be expressed as
Since the variation of strain is arbitrary, the stress–strain relation of an elastic material is given by
For a linear elastic material, the quantity is a linear function of , and can therefore be expressed as
where c is a fourth-rank tensor of material constants, also called the stiffness tensor. We can see why c must be a fourth-rank tensor by noting that, for a linear elastic material,
In index notation
The right-hand side constant requires four indices and is a fourth-rank quantity. We can also see that this quantity must be a tensor because it is a linear transformation that takes the strain tensor to the stress tensor. We can also show that the constant obeys the tensor transformation rules for fourth-rank tensors.
| Physical sciences | Solid mechanics | null |
229619 | https://en.wikipedia.org/wiki/Effusion | Effusion | In physics and chemistry, effusion is the process in which a gas escapes from a container through a hole of diameter considerably smaller than the mean free path of the molecules. Such a hole is often described as a pinhole and the escape of the gas is due to the pressure difference between the container and the exterior.
Under these conditions, essentially all molecules which arrive at the hole continue and pass through the hole, since collisions between molecules in the region of the hole are negligible. Conversely, when the diameter is larger than the mean free path of the gas, flow obeys the Sampson flow law.
In medical terminology, an effusion refers to accumulation of fluid in an anatomic space, usually without loculation. Specific examples include subdural, mastoid, pericardial and pleural effusions.
Etymology
The word effusion derives from the Latin word, effundo, which means "shed", "pour forth", "pour out", "utter", "lavish", "waste".
Into a vacuum
Effusion from an equilibrated container into outside vacuum can be calculated based on kinetic theory. The number of atomic or molecular collisions with a wall of a container per unit area per unit time (impingement rate) is given by:
assuming mean free path is much greater than pinhole diameter and the gas can be treated as an ideal gas.
If a small area on the container is punched to become a small hole, the effusive flow rate will be
where is the molar mass, is the Avogadro constant, and is the molar gas constant.
The average velocity of effused particles is
Combined with the effusive flow rate, the recoil/thrust force on the system itself is
An example is the recoil force on a balloon with a small hole flying in vacuum.
Measures of flow rate
According to the kinetic theory of gases, the kinetic energy for a gas at a temperature is
where is the mass of one molecule, is the root-mean-square speed of the molecules, and is the Boltzmann constant. The average molecular speed can be calculated from the Maxwell speed distribution
as (or, equivalently, ). The rate at which a gas of molar mass effuses (typically expressed as the number of molecules passing through the hole per second) is then
Here is the gas pressure difference across the barrier, is the area of the hole, is the Avogadro constant, is the gas constant and is the absolute temperature. Assuming the pressure difference between the two sides of the barrier is much smaller than , the average absolute pressure in the system (i.e. ), it is possible to express effusion flow as a volumetric flow rate as follows:
or
where is the volumetric flow rate of the gas, is the average pressure on either side of the orifice, and is the hole diameter.
Effect of molecular weight
At constant pressure and temperature, the root-mean-square speed and therefore the effusion rate are inversely proportional to the square root of the molecular weight. Gases with a lower molecular weight effuse more rapidly than gases with a higher molecular weight, so that the number of lighter molecules passing through the hole per unit time is greater.
Graham's law
Scottish chemist Thomas Graham (1805–1869) found experimentally that the rate of effusion of a gas is inversely proportional to the square root of the mass of its particles. In other words, the ratio of the rates of effusion of two gases at the same temperature and pressure is given by the inverse ratio of the square roots of the masses of the gas particles.
where and represent the molar masses of the gases.
This equation is known as Graham's law of effusion.
The effusion rate for a gas depends directly on the average velocity of its particles. Thus, the faster the gas particles are moving, the more likely they are to pass through the effusion orifice.
Knudsen cell
The Knudsen cell is used to measure the vapor pressures of a solid with very low vapor pressure. Such a solid forms a vapor at low pressure by sublimation. The vapor slowly effuses through a pinhole, and the loss of mass is proportional to the vapor pressure and can be used to determine this pressure. The heat of sublimation can also be determined by measuring the vapor pressure as a function of temperature, using the Clausius–Clapeyron relation.
| Physical sciences | Thermodynamics | Physics |
229643 | https://en.wikipedia.org/wiki/Molality | Molality | In chemistry, molality is a measure of the amount of solute in a solution relative to a given mass of solvent. This contrasts with the definition of molarity which is based on a given volume of solution.
A commonly used unit for molality is the moles per kilogram (mol/kg). A solution of concentration 1 mol/kg is also sometimes denoted as 1 molal. The unit mol/kg requires that molar mass be expressed in kg/mol, instead of the usual g/mol or kg/kmol.
Definition
The molality (b), of a solution is defined as the amount of substance (in moles) of solute, nsolute, divided by the mass (in kg) of the solvent, msolvent:
.
In the case of solutions with more than one solvent, molality can be defined for the mixed solvent considered as a pure pseudo-solvent. Instead of mole solute per kilogram solvent as in the binary case, units are defined as mole solute per kilogram mixed solvent.
Origin
The term molality is formed in analogy to molarity which is the molar concentration of a solution. The earliest known use of the intensive property molality and of its adjectival unit, the now-deprecated molal, appears to have been published by G. N. Lewis and M. Randall in the 1923 publication of Thermodynamics and the Free Energies of Chemical Substances. Though the two terms are subject to being confused with one another, the molality and molarity of a dilute aqueous solution are nearly the same, as one kilogram of water (solvent) occupies the volume of 1 liter at room temperature and a small amount of solute has little effect on the volume.
Unit
The SI unit for molality is moles per kilogram of solvent.
A solution with a molality of 3 mol/kg is often described as "3 molal", "3 m" or "3 m". However, following the SI system of units, the National Institute of Standards and Technology, the United States authority on measurement, considers the term "molal" and the unit symbol "m" to be obsolete, and suggests mol/kg or a related unit of the SI.
Usage considerations
Advantages
The primary advantage of using molality as a measure of concentration is that molality only depends on the masses of solute and solvent, which are unaffected by variations in temperature and pressure. In contrast, solutions prepared volumetrically (e.g. molar concentration or mass concentration) are likely to change as temperature and pressure change. In many applications, this is a significant advantage because the mass, or the amount, of a substance is often more important than its volume (e.g. in a limiting reagent problem).
Another advantage of molality is the fact that the molality of one solute in a solution is independent of the presence or absence of other solutes.
Problem areas
Unlike all the other compositional properties listed in "Relation" section (below), molality depends on the choice of the substance to be called “solvent” in an arbitrary mixture. If there is only one pure liquid substance in a mixture, the choice is clear, but not all solutions are this clear-cut: in an alcohol–water solution, either one could be called the solvent; in an alloy, or solid solution, there is no clear choice and all constituents may be treated alike. In such situations, mass or mole fraction is the preferred compositional specification.
Relation to other compositional quantities
In what follows, the solvent may be given the same treatment as the other constituents of the solution, such that the molality of the solvent of an n-solute solution, say b0, is found to be nothing more than the reciprocal of its molar mass, M0 (expressed in the unit kg/mol):
.
For the solutes the expression of molalities is similar:
.
The expressions linking molalities to mass fractions and mass concentrations contain the molar masses of the solutes Mi:
.
Similarly the equalities below are obtained from the definitions of the molalities and of the other compositional quantities.
The mole fraction of solvent can be obtained from the definition by dividing the numerator and denominator to the amount of solvent n0:
.
Then the sum of ratios of the other mole amounts to the amount of solvent is substituted with expressions from below containing molalities:
giving the result
.
Mass fraction
The conversions to and from the mass fraction, w1, of the solute in a single-solute solution are
where b1 is the molality and M1 is the molar mass of the solute.
More generally, for an n-solute/one-solvent solution, letting bi and wi be, respectively, the molality and mass fraction of the i-th solute,
,
where Mi is the molar mass of the ith solute, and w0 is the mass fraction of the solvent, which is expressible both as a function of the molalities as well as a function of the other mass fractions,
.
Substitution gives:
.
Mole fraction
The conversions to and from the mole fraction, x1 mole fraction of the solute in a single-solute solution are
,
where M0 is the molar mass of the solvent.
More generally, for an n-solute/one-solvent solution, letting xi be the mole fraction of the ith solute,
,
where x0 is the mole fraction of the solvent, expressible both as a function of the molalities as well as a function of the other mole fractions:
.
Substitution gives:
.
Molar concentration (molarity)
The conversions to and from the molar concentration, c1, for one-solute solutions are
,
where ρ is the mass density of the solution, b1 is the molality, and M1 is the molar mass (in kg/mol) of the solute.
For solutions with n solutes, the conversions are
,
where the molar concentration of the solvent c0 is expressible both as a function of the molalities as well as a function of the other molarities:
.
Substitution gives:
,
Mass concentration
The conversions to and from the mass concentration, ρsolute, of a single-solute solution are
,
or
,
where ρ is the mass density of the solution, b1 is the molality, and M1 is the molar mass of the solute.
For the general n-solute solution, the mass concentration of the ith solute, ρi, is related to its molality, bi, as follows:
,
where the mass concentration of the solvent, ρ0, is expressible both as a function of the molalities as well as a function of the other mass concentrations:
.
Substitution gives:
.
Equal ratios
Alternatively, one may use just the last two equations given for the compositional property of the solvent in each of the preceding sections, together with the relationships given below, to derive the remainder of properties in that set:
,
where i and j are subscripts representing all the constituents, the n solutes plus the solvent.
Example of conversion
An acid mixture consists of 0.76, 0.04, and 0.20 mass fractions of 70% HNO3, 49% HF, and H2O, where the percentages refer to mass fractions of the bottled acids carrying a balance of H2O. The first step is determining the mass fractions of the constituents:
.
The approximate molar masses in kg/mol are
.
First derive the molality of the solvent, in mol/kg,
,
and use that to derive all the others by use of the equal ratios:
.
Actually, bH2O cancels out, because it is not needed. In this case, there is a more direct equation: we use it to derive the molality of HF:
.
The mole fractions may be derived from this result:
,
,
.
Osmolality
Osmolality is a variation of molality that takes into account only solutes that contribute to a solution's osmotic pressure. It is measured in osmoles of the solute per kilogram of water. This unit is frequently used in medical laboratory results in place of osmolarity, because it can be measured simply by depression of the freezing point of a solution, or cryoscopy (see also: osmostat and colligative properties).
Relation to apparent (molar) properties
Molality appears in the expression of the apparent (molar) volume of a solute as a function of the molality b of that solute (and density of the solution and solvent):
,
.
For multicomponent systems the relation is slightly modified by the sum of molalities of solutes. Also a total molality and a mean apparent molar volume can be defined for the solutes together and also a mean molar mass of the solutes as if they were a single solute. In this case the first equality from above is modified with the mean molar mass M of the pseudosolute instead of the molar mass of the single solute:
,
, yi,j being ratios involving molalities of solutes i,j and the total molality bT.
The sum of products molalities - apparent molar volumes of solutes in their binary solutions equals the product between the sum of molalities of solutes and apparent molar volume in ternary or multicomponent solution.
.
Relation to apparent molar properties and activity coefficients
For concentrated ionic solutions the activity coefficient of the electrolyte is split into electric and statistical components.
The statistical part includes molality b, hydration index number h, the number of ions from the dissociation and the ratio ra between the apparent molar volume of the electrolyte and the molar volume of water.
Concentrated solution statistical part of the activity coefficient is:
.
Molalities of a ternary or multicomponent solution
The molalities of solutes b1, b2 in a ternary solution obtained by mixing two binary aqueous solutions with different solutes (say a sugar and a salt or two different salts) are different than the initial molalities of the solutes bii in their binary solutions:
,
,
,
.
The content of solvent in mass fractions w01 and w02 from each solution of masses ms1 and ms2 to be mixed as a function of initial molalities is calculated. Then the amount (mol) of solute from each binary solution is divided by the sum of masses of water after mixing:
,
.
Mass fractions of each solute in the initial solutions w11 and w22
are expressed as a function of the initial molalities b11, b22:
,
.
These expressions of mass fractions are substituted in the final molalitaties:
,
.
The results for a ternary solution can be extended to a multicomponent solution (with more than two solutes).
From the molalities of the binary solutions
The molalities of the solutes in a ternary solution can be expressed also from molalities in the binary solutions and their masses:
,
.
The binary solution molalities are:
,
.
The masses of the solutes determined from the molalities of the solutes and the masses of water can be substituted in the expressions of the masses of solutions:
.
Similarly for the mass of the second solution:
.
One can obtain the masses of water present in the sum from the denominator of the molalities of the solutes in the ternary solutions as functions of binary molalities and masses of solution:
,
.
Thus the ternary molalities are:
,
.
For solutions with three or more solutes the denominator is a sum of the masses of solvent in the n binary solutions which are mixed:
,
,
.
| Physical sciences | Concentration | Basics and measurement |
229714 | https://en.wikipedia.org/wiki/Carriage | Carriage | A carriage is a two- or four-wheeled horse-drawn vehicle for passengers. In Europe they were a common mode of transport for the wealthy during the Roman Empire, and then again from around 1600 until they were replaced by the motor car around 1900. They were generally owned by the rich, but second-hand private carriages became common public transport, the equivalent of modern cars used as taxis. Carriage suspensions are by leather strapping or, on those made in recent centuries, steel springs. There are numerous names for different types. Two-wheeled carriages are usually owner-driven.
Coaches are a special category within carriages. They are carriages with four corner posts and a fixed roof. Two-wheeled war chariots and transport vehicles such as four-wheeled wagons and two-wheeled carts were forerunners of carriages.
In the 21st century, horse-drawn carriages are occasionally used for public parades by royalty and for traditional formal ceremonies. Simplified modern versions are made for tourist transport in warm countries and for those cities where tourists expect open horse-drawn carriages to be provided. Simple metal sporting versions are still made for the sport known as competitive driving.
Overview
The word carriage (abbreviated carr or cge) is from Old Northern French , to carry in a vehicle. The word car, then meaning a kind of two-wheeled cart for goods, also came from Old Northern French about the beginning of the 14th century (probably derived from the Late Latin , a car); it is also used for railway carriages and in the US around the end of the 19th century, early cars (automobiles) were briefly called horseless carriages.
History
Early history
Some horse carts found in Celtic graves show hints that their platforms were suspended elastically. Four-wheeled wagons were used in Bronze Age Europe, and their form known from excavations suggests that the basic construction techniques of wheel and undercarriage (that survived until the age of the motor car) were established then.
First prototyped in the 3rd millennium BC, a bullock cart is a large two-wheeled cart pulled by oxen or buffalo. It includes a sturdy wooden pole between the oxen, a yoke connecting a pair of oxen, a wooden platform for passengers or cargo, and large steel rimmed wooden wheels.
Two-wheeled carriage models have been discovered from the Indus valley civilization including twin horse drawn covered carriages resembling ekka from various sites such as Harappa, Mohenjo Daro and Chanhu Daro. The earliest recorded sort of carriage was the chariot, reaching Mesopotamia as early as 1900 BC. Used typically for warfare by Egyptians, the Near Easterners and Europeans, it was essentially a two-wheeled light basin carrying one or two standing passengers, drawn by one to two horses. The chariot was revolutionary and effective because it delivered fresh warriors to crucial areas of battle with swiftness.
Roman carriage
First century BC Romans used sprung wagons for overland journeys. It is likely that Roman carriages employed some form of suspension on chains or leather straps, as indicated by carriage parts found in excavations. In 2021 archaeologists discovered the remains of a ceremonial four wheel carriage, a pilentum, near the ancient Roman city of Pompeii. It is thought the pilentum may have been used in ceremonies such as weddings. The find has been described as being "in an excellent state of preservation".
Ancient Chinese carriage
Though the exact date of when the Chinese started to use carriages is largely unknown, early oracle bone inscriptions discovered in Henan province show that the carriage had already developed into many different forms.
The earliest archaeological evidence of chariots in China, a chariot burial site discovered in 1933 at Hougang, Anyang in Henan province, dates to the rule of King Wu Ding of the late Shang dynasty (). Oracle bone inscriptions suggest that the western enemies of the Shang used limited numbers of chariots in battle, but the Shang themselves used them only as mobile command-vehicles and in royal hunts.
During the Shang dynasty, members of the royal family were buried with a complete household and servants, including a chariot, horses, and a charioteer. A Shang chariot was often drawn by two horses, but four-horse variants are occasionally found in burials.
Jacques Gernet claims that the Zhou dynasty, which conquered the Shang ca. 1046 BCE, made more use of the chariot than did the Shang and "invented a new kind of harness with four horses abreast". The crew consisted of an archer, a driver, and sometimes a third warrior who was armed with a spear or dagger-axe. From the 8th to 5th centuries BCE the Chinese use of chariots reached its peak. Although chariots appeared in greater numbers, infantry often defeated charioteers in battle.
Massed-chariot warfare became all but obsolete after the Warring-States Period (476–221 BCE). The main reasons were increased use of the crossbow, use of long halberds up to long and pikes up to long, and the adoption of standard cavalry units, and the adaptation of mounted archery from nomadic cavalry, which were more effective. Chariots would continue to serve as command posts for officers during the Qin dynasty (221–206 BCE) and the Han dynasty (206 BCE–220 CE), while armored chariots were also used during the Han dynasty against the Xiongnu Confederation in the Han–Xiongnu War (133 BC to 89 CE), specifically at the Battle of Mobei (119 BCE).
Before the Han dynasty, the power of Chinese states and dynasties was often measured by the number of chariots they were known to have. A country of a thousand chariots ranked as a medium country, and a country of ten thousand chariots ranked as a huge and powerful country.
Medieval carriage
The medieval carriage was typically a four-wheeled wagon type, with a rounded top ("tilt") similar in appearance to the Conestoga Wagon familiar from the United States. Sharing the traditional form of wheels and undercarriage known since the Bronze Age, it very likely also employed the pivoting fore-axle in continuity from the ancient world. Suspension (on chains) is recorded in visual images and written accounts from the 14th century ("chars branlant" or rocking carriages), and was in widespread use by the 15th century. Carriages were largely used by royalty, aristocrats (and especially by women), and could be elaborately decorated and gilded. These carriages were usually on four wheels and were drawn by two to four horses depending on their size and status. Wood and iron were the primary materials needed to build a carriage and carriages that were used by non-royalty were covered by plain leather.
Another form of carriage was the pageant wagon of the 14th century. Historians debate the structure and size of pageant wagons; however, they are generally miniature house-like structures that rest on four to six wheels depending on the size of the wagon. The pageant wagon is significant because up until the 14th century most carriages were on two or three wheels; the chariot, rocking carriage, and baby carriage are two examples of carriages which pre-date the pageant wagon. Historians also debate whether or not pageant wagons were built with pivotal axle systems, which allowed the wheels to turn. Whether it was a four- or six-wheel pageant wagon, most historians maintain that pivotal axle systems were implemented on pageant wagons because many roads were often winding with some sharp turns. Six wheel pageant wagons also represent another innovation in carriages; they were one of the first carriages to use multiple pivotal axles. Pivotal axles were used on the front set of wheels and the middle set of wheels. This allowed the horse to move freely and steer the carriage in accordance with the road or path.
Coach
One of the great innovations in carriage history was the invention of the suspended carriage or the chariot branlant (though whether this was a Roman or medieval innovation remains uncertain). The "chariot branlant" of medieval illustrations was suspended by chains rather than leather straps as had been believed. Suspension, whether on chains or leather, might provide a smoother ride since the carriage body no longer rested on the axles, but could not prevent swinging (branlant) in all directions. It is clear from illustrations (and surviving examples) that the medieval suspended carriage with a round tilt was a widespread European type, referred to by any number of names (car, currus, char, chariot).
In 14th century England carriages, like the one illustrated in the Luttrell Psalter, would still have been a quite rare means of aristocratic transport, and they would have been very costly until the end of the century. They would have had four six-spoke six-foot high wheels that were linked by greased axles under the body of the coach, and did not necessarily have any suspension. The chassis was made from oak beam and the barrel shaped roof was covered in brightly painted leather or cloth. The interior would include seats, beds, cushions, tapestries and even rugs. They would be pulled by four to five horses.
Under King Mathias Corvinus (1458–90), who enjoyed fast travel, the Hungarians developed fast road transport, and the town of Kocs between Budapest and Vienna became an important post-town, and gave its name to the new vehicle type. The earliest illustrations of the Hungarian "Kochi-wagon" do not indicate any suspension, a body with high sides of lightweight wickerwork, and typically drawn by three horses in harness. Later models were considerably lighter and famous for a single horse being able to draw many passengers.
The Hungarian coach spread across Europe, initially rather slowly, in part due to Ippolito d'Este of Ferrara (1479–1529), nephew of Mathias' queen Beatrix of Aragon, who as a very junior Archbishopric of Esztergom developed a taste for Hungarian riding and took his carriage and driver back to Italy. Then rather suddenly, in around 1550, the "coach" made its appearance throughout the major cities of Europe, and the new word entered the vocabulary of all their languages. However, the new "coach" seems to have been a fashionable concept (fast road travel for men) as much as any particular type of vehicle, and there is no obvious technological change that accompanied the innovation, either in the use of suspension (which came earlier), or the adoption of springs (which came later). As its use spread throughout Europe in the late 16th century, the coach's body structure was ultimately changed, from a round-topped tilt to the "four-poster" carriages that became standard everywhere by c.1600.
Later development of the coach
The coach had doors in the side, with an iron step protected by leather that became the "boot" in which servants might ride. The driver sat on a seat at the front, and the most important occupant sat in the back facing forwards. The earliest coaches can be seen at Veste Coburg, Lisbon, and the Moscow Kremlin, and they become a commonplace in European art. It was not until the 17th century that further innovations with steel springs and glazing took place, and only in the 18th century, with better road surfaces, was there a major innovation with the introduction of the steel C-spring.
Many innovations were proposed, and some patented, for new types of suspension or other features. It was only from the 18th century that changes to steering systems were suggested, including the use of the 'fifth wheel' substituted for the pivoting fore-axle, and on which the carriage turned. Another proposal came from Erasmus Darwin, a young English doctor who was driving a carriage about 10,000 miles a year to visit patients all over England. Darwin found two essential problems or shortcomings of the commonly used light carriage or Hungarian carriage. First, the front wheels were turned by a pivoting front axle, which had been used for years, but these wheels were often quite small and hence the rider, carriage and horse felt the brunt of every bump on the road. Secondly, he recognized the danger of overturning.
A pivoting front axle changes a carriage's base from a rectangle to a triangle because the wheel on the inside of the turn is able to turn more sharply than the outside front wheel. Darwin suggested a fix for these insufficiencies by proposing a principle in which the two front wheels turn (independently of the front axle) about a centre that lies on the extended line of the back axle. This idea was later patented in 1818 as Ackermann steering. Darwin argued that carriages would then be easier to pull and less likely to overturn.
Carriage use in North America came with the establishment of European settlers. Early colonial horse tracks quickly grew into roads especially as the colonists extended their territories southwest. Colonists began using carts as these roads and trading increased between the north and south. Eventually, carriages or coaches were sought to transport goods as well as people. As in Europe, chariots, coaches and/or carriages were a mark of status. The tobacco planters of the South were some of the first Americans to use the carriage as a form of human transportation. As the tobacco farming industry grew in the southern colonies so did the frequency of carriages, coaches and wagons. Upon the turn of the 18th century, wheeled vehicle use in the colonies was at an all-time high. Carriages, coaches and wagons were being taxed based on the number of wheels they had. These taxes were implemented in the South primarily as the South had superior numbers of horses and wheeled vehicles when compared to the North. Europe, however, still used carriage transportation far more often and on a much larger scale than anywhere else in the world.
Demise
Carriages and coaches began to disappear as use of steam propulsion began to generate more and more interest and research. Steam power quickly won the battle against animal power as is evident by a newspaper article written in England in 1895 entitled "Horseflesh vs. Steam". The article highlights the death of the carriage as the main means of transportation.
Today
Today, carriages are still used for day-to-day transport in the United States by some minority groups such as the Amish. They are also still used in tourism as vehicles for sightseeing in cities such as Bruges, Vienna, New Orleans, and Little Rock, Arkansas.
The most complete working collection of carriages can be seen at the Royal Mews in London where a large selection of vehicles is in regular use. These are supported by a staff of liveried coachmen, footmen and postillions. The horses earn their keep by supporting the work of the Royal Household, particularly during ceremonial events. Horses pulling a large carriage known as a "covered brake" collect the Yeoman of the Guard in their distinctive red uniforms from St James's Palace for Investitures at Buckingham Palace; High Commissioners or Ambassadors are driven to their audiences with the King and Queen in landaus; visiting heads of state are transported to and from official arrival ceremonies and members of the Royal Family are driven in Royal Mews coaches during Trooping the Colour, the Order of the Garter service at Windsor Castle and carriage processions at the beginning of each day of Royal Ascot.
Construction
Body
Carriages may be enclosed or open, depending on the type. The top cover for the body of a carriage, called the head or hood, is often flexible and designed to be folded back when desired. Such a folding top is called a bellows top or calash. A hoopstick forms a light framing member for this kind of hood. The top, roof or second-story compartment of a closed carriage, especially a diligence, was called an imperial. A closed carriage may have side windows called quarter lights (British) as well as windows in the doors, hence a "glass coach". On the forepart of an open carriage, a screen of wood or leather called a dashboard intercepts water, mud or snow thrown up by the heels of the horses. The dashboard or carriage top sometimes has a projecting sidepiece called a wing (British). A foot iron or footplate may serve as a carriage step.
A carriage driver sits on a box or perch, usually elevated and small. When at the front, it is known as a dickey box, a term also used for a seat at the back for servants. A footman might use a small platform at the rear called a footboard or a seat called a rumble behind the body. Some carriages have a moveable seat called a jump seat. Some seats had an attached backrest called a lazyback.
The shafts of a carriage were called limbers in English dialect. Lancewood, a tough elastic wood of various trees, was often used especially for carriage shafts. A holdback, consisting of an iron catch on the shaft with a looped strap, enables a horse to back or hold back the vehicle. The end of the tongue of a carriage is suspended from the collars of the harness by a bar called the yoke. At the end of a trace, a loop called a cockeye attaches to the carriage.
In some carriage types, the body is suspended from several leather straps called braces or thoroughbraces, attached to or serving as springs.
Undercarriage
Beneath the carriage body is the undergear or undercarriage (or simply carriage), consisting of the running gear and chassis. The wheels and axles, in distinction from the body, are the running gear. The wheels revolve upon bearings or a spindle at the ends of a bar or beam called an axle or axletree. Most carriages have either one or two axles. On a four-wheeled vehicle, the forward part of the running gear, or forecarriage, is arranged to permit the front axle to turn independently of the fixed rear axle. In some carriages a dropped axle, bent twice at a right angle near the ends, allows for a low body with large wheels. A guard called a dirtboard keeps dirt from the axle arm.
Several structural members form parts of the chassis supporting the carriage body. The fore axletree and the splinter bar above it (supporting the springs) are united by a piece of wood or metal called a futchel, which forms a socket for the pole that extends from the front axle. For strength and support, a rod called the backstay may extend from either end of the rear axle to the reach, the pole or rod joining the hind axle to the forward bolster above the front axle.
A skid called a drag, dragshoe, shoe or skidpan retards the motion of the wheels. A London patent of 1841 describes one such apparatus: "An iron-shod beam, slightly longer than the radius of the wheel, is hinged under the axle so that when it is released to strike the ground the forward momentum of the vehicle wedges it against the axle". The original feature of this modification was that instead of the usual practice of having to stop the carriage to retract the beam and so lose useful momentum the chain holding it in place is released (from the driver's position) so that it is allowed to rotate further in its backwards direction, releasing the axle. A system of "pendant-levers" and straps then allows the beam to return to its first position and be ready for further use.
A catch or block called a trigger may be used to hold a wheel on an incline.
A horizontal wheel or segment of a wheel called a fifth wheel sometimes forms an extended support to prevent the carriage from tipping; it consists of two parts rotating on each other about the kingbolt or perchbolt above the fore axle and beneath the body. A block of wood called a headblock might be placed between the fifth wheel and the forward spring.
Fittings, furnishings and appointments
Originally, the word fittings referred to metal elements such as bolts and brackets, furnishings leaned more to leatherwork and upholstery or referred to metal buckles on harness, and appointments were things brought to a carriage but not part of it, however all of these words have blended together over time and are often used interchangeably to mean the smaller components or parts of a carriage or equipment. All the shiny metal fittings on a vehicle should be one color, such as brass (yellow) or nickel (white), and should match the buckle color of any harness used with the vehicle. Early bodies of horseless carriages were constructed by coachmakers using the same parts used in carriages and coaches, and some horse carriage terminology has survived in modern automobiles.
Upholstery: Seats might be upholstered using leather, broadcloth, or plush fabrics. Elegant carriages might have upholstery-lined walls and ceilings, and button-tucked velvet seats trimmed with gold braid.
Carriage lamps: First used around 1700, oil-powered lamps were used throughout the 1800s, though abandoned in favor of candles in the late 1800s, as oil was messy. Lamps are mounted on lamp brackets and are removable for storage, daily wick trimming, or during daylight hours.
Boot: Any of several box-like parts of a carriage used for storage of small items. A boot may be found under the coachman's seat, under the passenger's seat, or behind the body of the carriage between the rear wheels. This led to the use of the term boot in British English for the main storage compartment of an automobile.
Whip socket: Tubular holder for a whip usually mounted on the dashboard or to the right of the driver.
Whip: A long whip composed of a stiff stick (called the stock), a long flexible thong, and a short lash. The length should be appropriate for the distance from the driver (who is also called a Whip) to the shoulder of the forwardmost horse. With a small pony and cart a whip of overall length of 7 or 8 feet might be appropriate, whereas driving a team of four horses might require an overall length of 17 feet. Driving whips are not "cracked" to make noise, but are a communication aid used by touching the lash on or near the shoulder of the horse.
Blankets: in cold weather, blankets for the driver and passengers and often horse blankets as well may be carried in a boot.
Carriage terminology
The carriage driver is called a whip. A person whose business was to drive a carriage was a coachman. A person dressed in livery is called a footman. An attendant on horseback called an outrider. A carriage starter directed the flow of vehicles taking on passengers at the curbside. A hackneyman hired out horses and carriages.
Upper-class people of wealth and social position, those wealthy enough to keep carriages, were referred to as carriage folk or carriage trade.
Carriage passengers often used a lap robe as a blanket or similar covering for their legs, lap and feet.
A horse especially bred for carriage use by appearance and stylish action is called a carriage horse; one for use on a road is a road horse. One such breed is the Cleveland Bay, uniformly bay in color, of good conformation and strong constitution. Horses were broken in using a bodiless carriage frame called a break or brake.
A carriage dog or coach dog is bred for running beside a carriage.
A roofed structure that extends from the entrance of a building over an adjacent driveway and that shelters callers as they get in or out of their vehicles is known as a carriage porch or porte cochere. An outbuilding for a carriage is a coach house, which was often combined with accommodation for a groom or other servants.
A livery stable kept horses and usually carriages for hire. A range of stables, usually with carriage houses (remises) and living quarters built around a yard, court or street, is called a mews.
A kind of dynamometer called a peirameter indicates the power necessary to haul a carriage over a road or track.
Competitive driving
In most European and English-speaking countries, driving is a competitive equestrian sport. Many horse shows host driving competitions for a particular style of driving, breed of horse, or type of vehicle. Show vehicles are usually carriages, carts, or buggies and, occasionally, sulkies or wagons. Modern high-technology carriages are made purely for competition by companies such as Bennington Carriages. in England.
Terminology varies: the simple, lightweight two- or four-wheeled show vehicle common in many nations is called a "cart" in the US, but a "carriage" in Australia.
Internationally, there is intense competition in the all-round test of driving: combined driving, also known as horse-driving trials, an equestrian discipline regulated by the Fédération Équestre Internationale (International Equestrian Federation) with national organizations representing each member country. World championships are conducted in alternate years, including single-horse, horse pairs and four-in-hand championships. The World Equestrian Games, held at four-year intervals, also includes a four-in-hand competition.
For pony drivers, the World Combined Pony Championships are held every two years and include singles, pairs and four-in-hand events.
Carriage museums and collections
Argentina
Muhfit (Museo Histórico Fuerte Independencia Tandil), Tandil.
Australia
Cobb & Co Museum – National Carriage Collection, Queensland Museum, Toowoomba, Queensland.
National Trust of Australia (Victoria) Carriage Collection
Austria
Imperial Carriage Museum at Schönbrunn Palace in Vienna
Kutschenmuseum in Laa an der Thaya
Belgium
Bornem Castle Carriage collection in Bornem
in Bree
De Groom Carriage Center Bruges in Bruges
Koetsen Verdonckt in Maarkedal
Royal Museum of Art and History in Brussels
Brazil
Imperial Museum in Petrópolis
National Historical Museum in Rio de Janeiro
Canada
Campbell Carriage Factory Museum in Sackville, New Brunswick
Kings Landing Historical Settlement in Prince William, New Brunswick — large collection of horse and oxen drawn vehicles
Remington Carriage Museum in Cardston, Alberta
Denmark
Royal Carriage Museum, Christiansborg Palace in Copenhagen
Egypt
Carriage Museum in Cairo Citadel
France
Apremont-sur-Allier, Musée des calèches (Berry)
Bourg, Musée Au temps des calèches (Guyenne)
Cazes-Mondenard, Musée de l'Attelage et du corbillard Yvan Quercy (Quercy)
Château de Chambord. Carriage room of the Count of Chambord in Chambord, Loir-et-Cher
Cussac-Fort-Médoc, Musée du cheval du château Lanessan (Guyenne)
Le Fleix, Musée de l’hippomobile André Clament (Périgord)
Les Épesses, Musée de la voiture à cheval (Vendée, Bas-Poitou)
Marcigny, Musée de la voiture à cheval (Bourgogne)
National Car and Tourism Museum at Château de Compiègne in Compiègne
Plouay. Musée du conservatoire de la voiture hippomobile (Brittany)
Sacy-le-grand, Musée du cheval de trait (Picardie)
Saint-Auvent, musée Au temps jadis (Limousin)
Sérignan, Musée de l’attelage et du cheval (Languedoc)
Château de Vaux-le-Vicomte, Musée des Equipages (Île-de-France)
Galerie des Carrosses at Grande Écurie in Versailles (Île-de-France)
Germany
Hesse Museum of Carriages and Sleighs in Lohfelden near Kassel
Marstallmuseum of Carriages and Sleighs in the former Royal Stables, Nymphenburg Palace, Munich
Romano-Germanic Museum
Italy
Collection at CastelBrando near Cison di Valmarino
Museo "Le Carrozze d'Epoca", Rome.
Museo Civico delle Carrozze d'Epoca di Codroipo.
Museo Civico delle Carrozze d'Epoca, San Martino, Udine.
Museo della Carrozza in Macerata
Museo delle Carrozze del Quirinale, Rome.
Museum of Coaches at Palazzo Farnese, Piacenza in Piacenza
Carriage exhibit of the Grand Ducal court at Palazzo Pitti in Florence
Museo delle Carrozze, Catanzaro.
Carriage collection at Villa Barbaro in Maser, Veneto
Carriage collection at Villa Pignatelli in Naples
Japan
Japanese Imperial Household Agency, Tokyo
Netherlands
, Leek in Groningen.
Poland
Kozłówka Palace in Kozłówka
Łańcut Castle in Łańcut
in Rogalin
Portugal
Geraz do Lima Carriage museum in Viana do Castelo
National Coach Museum () in Lisbon
Spain
, Seville
Igualada Muleteer's Museum in Igualada
Sweden
Ulriksdal Palace in Edsviken
Switzerland
Basel Historical Museum in Basel
Turkey
Tofaş Museum of Cars and Anatolian Carriages in Bursa
United Kingdom
Alnwick Castle in Alnwick, Northumberland
Arlington Court & The National Trust's Carriage Collection in Arlington, Devon
Balmoral Castle in Aberdeenshire, Scotland
Gordon Boswell Romany Museum in Spalding, Lincolnshire
Mossman Carriage Collection in Luton, Bedfordshire
Royal Mews at Buckingham Palace in London
Sandringham House in Sandringham, Norfolk
Swingletree Carriage Collection of John Parker in Diss, Norfolk
Tyrwhitt-Drake Museum of Carriages in Maidstone, Kent
United States
Angels Camp Museum in Angels Camp, California
Carriage Museum of America, Lexington, Kentucky
Florida Carriage Museum & Resort in Weirsdale, Florida (formerly Austin Carriage Museum)
Forney Transportation Museum in Denver, Colorado
Frick Car & Carriage Museum in Pittsburgh, Pennsylvania, preserving carriages owned by Henry Clay Frick and his family
Genesee Country Village and Museum in Wheatland, New York
Granger Homestead and Carriage Museum in Canandaigua, New York
Harness Racing Museum & Hall of Fame in Goshen, New York
Henry Ford Museum in Dearborn, Michigan
Horseshoe Barn and Annex at Shelburne Museum in Shelburne, Vermont
Jeremiah Reeves House and Carriage House in Dover, Ohio
Long Island Museum of American Art, History, and Carriages in Stony Brook, New York
Maymont in Richmond, Virginia
Morven Park's Winmill Carriage Museum in Leesburg, Virginia
Northwest Carriage Museum in Raymond, Washington
Pioneer Village in Farmington, Utah
Robert H. Renneberger Carriage Museum in Frederick, Maryland
Robert Thomas Carriage Museum in Blackstone, Virginia
Skyline Farm Carriage Museum, North Yarmouth, Maine
Thrasher Carriage Collection at Allegany Museum in Cumberland, Maryland
Washington, Kentucky Carriage Museum
Wesley Jung Carriage Museum on Wade House Historic Site in Greenbush, Wisconsin
William A. Heiss House and Buggy Shop in Mifflinburg, Pennsylvania; includes 19th century carriage factory
Types of horse-drawn carriages
Numerous varieties of horse-drawn carriages existed, Arthur Ingram's Horse Drawn Vehicles since 1760 in Colour lists 325 types with a short description of each. By the early 19th century one's choice of carriage was only in part based on practicality and performance; it was also a status statement and subject to changing fashions.
| Technology | Animal-powered transport | null |
229914 | https://en.wikipedia.org/wiki/Ape | Ape | Apes (collectively Hominoidea ) are a clade of Old World simians native to sub-Saharan Africa and Southeast Asia (though they were more widespread in Africa, most of Asia, and Europe in prehistory, and counting humans are found globally). Apes are more closely related to Old World monkeys (family Cercopithecidae) than to the New World monkeys (Platyrrhini) with both Old World monkeys and apes placed in the clade Catarrhini. Apes do not have tails due to a mutation of the TBXT gene. In traditional and non-scientific use, the term ape can include tailless primates taxonomically considered Cercopithecidae (such as the Barbary ape and black ape), and is thus not equivalent to the scientific taxon Hominoidea. There are two extant branches of the superfamily Hominoidea: the gibbons, or lesser apes; and the hominids, or great apes.
The family Hylobatidae, the lesser apes, include four genera and a total of 20 species of gibbon, including the lar gibbon and the siamang, all native to Asia. They are highly arboreal and bipedal on the ground. They have lighter bodies and smaller social groups than great apes.
The family Hominidae (hominids), the great apes, include four genera comprising three extant species of orangutans and their subspecies, two extant species of gorillas and their subspecies, two extant species of chimpanzees and their subspecies, and humans in a single extant subspecies.
Except for gorillas and humans, hominoids are agile climbers of trees. Apes eat a variety of plant and animal foods, with the majority of food being plant foods, which can include fruits, leaves, stalks, roots and seeds, including nuts and grass seeds. Human diets are sometimes substantially different from that of other hominoids due in part to the development of technology and a wide range of habitation.
All non-human hominoids are rare and threatened with extinction. The main threat is habitat loss, though some populations are further imperiled by hunting. The great apes of Africa are also facing threat from the Ebola virus.
Name and terminology
"Ape", from Old English apa, is a word of uncertain origin. The term has a history of rather imprecise usage—and of comedic or punning usage in the vernacular. Its earliest meaning was generally of any non-human anthropoid primate, as is still the case for its cognates in other Germanic languages.
Later, after the term "monkey" had been introduced into English, "ape" was specialized to refer to a tailless (therefore exceptionally human-like) primate. Thus, the term "ape" obtained two different meanings, as shown in the 1911 Encyclopædia Britannica entry: it could be used as a synonym for "monkey" and it could denote the tailless human-like primate in particular.
Some, or recently all, hominoids are also called "apes", but the term is used broadly and has several different senses within both popular and scientific settings. "Ape" has been used as a synonym for "monkey" or for naming any primate with a human-like appearance, particularly those without a tail. Biologists have traditionally used the term "ape" to mean a member of the superfamily Hominoidea other than humans, but more recently to mean all members of Hominoidea. So "ape"—not to be confused with "great ape"—now becomes another word for hominoid including humans.
The taxonomic term hominoid is derived from, and intended as encompassing, the hominids, the family of great apes. Both terms were introduced by Gray (1825). The term hominins is also due to Gray (1824), intended as including the human lineage (see also Hominidae#Terminology, Human taxonomy).
The distinction between apes and monkeys is complicated by the traditional paraphyly of monkeys: Apes emerged as a sister group of Old World Monkeys in the catarrhines, which are a sister group of New World Monkeys. Therefore, cladistically, apes, catarrhines and related contemporary extinct groups such as Parapithecidae are monkeys as well, for any consistent definition of "monkey". "Old World monkey" may also legitimately be taken to be meant to include all the catarrhines, including apes and extinct species such as Aegyptopithecus, in which case the apes, Cercopithecoidea and Aegyptopithecus emerged within the Old World monkeys.
The primates called "apes" today became known to Europeans after the 18th century. As zoological knowledge developed, it became clear that taillessness occurred in a number of different and otherwise distantly related species. Sir Wilfrid Le Gros Clark was one of those primatologists who developed the idea that there were trends in primate evolution and that the extant members of the order could be arranged in an "ascending series", leading from "monkeys" to "apes" to humans. Within this tradition "ape" came to refer to all members of the superfamily Hominoidea except humans. As such, this use of "apes" represented a paraphyletic grouping, meaning that, even though all species of apes were descended from a common ancestor, this grouping did not include all the descendant species, because humans were excluded from being among the apes.
Traditionally, the English-language vernacular name "apes" does not include humans, but phylogenetically, humans (Homo) form part of the family Hominidae within Hominoidea. Thus, there are at least three common, or traditional, uses of the term "ape": non-specialists may not distinguish between "monkeys" and "apes", that is, they may use the two terms interchangeably; or they may use "ape" for any tailless monkey or non-human hominoid; or they may use the term "ape" to just mean the non-human hominoids.
Modern taxonomy aims for the use of monophyletic groups for taxonomic classification;
Some literature may now use the common name "ape" to mean all members of the superfamily Hominoidea, including humans. For example, in his 2005 book, Benton wrote "The apes, Hominoidea, today include the gibbons and orang-utan ... the gorilla and chimpanzee ... and humans". Modern biologists and primatologists refer to apes that are not human as "non-human" apes. Scientists broadly, other than paleoanthropologists, may use the term "hominin" to identify the human clade, replacing the term "hominid". See terminology of primate names.
See below, History of hominoid taxonomy, for a discussion of changes in scientific classification and terminology regarding hominoids.
Evolution
Although the hominoid fossil record is still incomplete and fragmentary, there is now enough evidence to provide an outline of the evolutionary history of humans. Previously, the divergence between humans and other extant hominoids was thought to have occurred 15 to 20 million years ago, and several species of that time period, such as Ramapithecus, were once thought to be hominins and possible ancestors of humans. But, later fossil finds indicated that Ramapithecus was more closely related to the orangutan; and new biochemical evidence indicates that the last common ancestor of humans and non-hominins (that is, the chimpanzees) occurred between 5 and 10 million years ago, and probably nearer the lower end of that range (more recent); see Chimpanzee–human last common ancestor (CHLCA).
Taxonomic classification and phylogeny
Genetic analysis combined with fossil evidence indicates that hominoids diverged from the Old World monkeys about 25 million years ago (mya), near the Oligocene–Miocene boundary. The gibbons split from the rest about 18 mya, and the hominid splits happened 14 mya (Pongo), 7 mya (Gorilla), and 3–5 mya (Homo & Pan). In 2015, a new genus and species were described, Pliobates cataloniae, which lived 11.6 mya, and appears to predate the split between Hominidae and Hylobatidae.
The families, and extant genera and species of hominoids are:
Superfamily Hominoidea
Family Hominidae: hominids ("great apes")
Genus Pongo: orangutans
Bornean orangutan, P. pygmaeus
Sumatran orangutan, P. abelii
Tapanuli orangutan, P. tapanuliensis
Genus Gorilla: gorillas
Western gorilla, G. gorilla
Eastern gorilla, G. beringei
Genus Homo: humans
Human, H. sapiens
Genus Pan: chimpanzees
Chimpanzee, P. troglodytes
Bonobo, P. paniscus
Family Hylobatidae: gibbons ("lesser apes")
Genus Hylobates
Lar gibbon or white-handed gibbon, H. lar
Bornean white-bearded gibbon, H. albibarbis
Agile gibbon or black-handed gibbon, H. agilis
Western grey gibbon or Abbott's grey gibbon, H. abbotti
Eastern grey gibbon or northern grey gibbon, H. funereus
Müller's gibbon or southern grey gibbon, H. muelleri
Silvery gibbon, H. moloch
Pileated gibbon or capped gibbon, H. pileatus
Kloss's gibbon or Mentawai gibbon or bilou, H. klossii
Genus Hoolock
Western hoolock gibbon, H. hoolock
Eastern hoolock gibbon, H. leuconedys
Skywalker hoolock gibbon, H. tianxing
Genus Symphalangus
Siamang, S. syndactylus
Genus Nomascus
Northern buffed-cheeked gibbon, N. annamensis
Black crested gibbon, N. concolor
Eastern black crested gibbon, N. nasutus
Hainan black crested gibbon, N. hainanus
Southern white-cheeked gibbon N. siki
White-cheeked crested gibbon, N. leucogenys
Yellow-cheeked gibbon, N. gabriellae
History of hominoid taxonomy
The history of hominoid taxonomy is complex and somewhat confusing. Recent evidence has changed our understanding of the relationships between the hominoids, especially regarding the human lineage; and the traditionally used terms have become somewhat confused. Competing approaches to methodology and terminology are found among current scientific sources. Over time, authorities have changed the names and the meanings of names of groups and subgroups as new evidence — that is, new discoveries of fossils and tools and of observations in the field, plus continual comparisons of anatomy and DNA sequences — has changed the understanding of relationships between hominoids. There has been a gradual demotion of humans from being 'special' in the taxonomy to being one branch among many. This recent turmoil (of history) illustrates the growing influence on all taxonomy of cladistics, the science of classifying living things strictly according to their lines of descent.
Today, there are eight extant genera of hominoids. They are the four genera in the family Hominidae, namely Homo, Pan, Gorilla, and Pongo; plus four genera in the family Hylobatidae (gibbons): Hylobates, Hoolock, Nomascus and Symphalangus. (The two subspecies of hoolock gibbons were recently moved from the genus Bunopithecus to the new genus Hoolock and re-ranked as species; a third species was described in January 2017).
In 1758, Carl Linnaeus, relying on second- or third-hand accounts, placed a second species in Homo along with H. sapiens: Homo troglodytes ("cave-dwelling man"). Although the term "Orang Outang" is listed as a variety – Homo sylvestris – under this species, it is nevertheless not clear to which animal this name refers, as Linnaeus had no specimen to refer to, hence no precise description. Linnaeus may have based Homo troglodytes on reports of mythical creatures, then-unidentified simians, or Asian natives dressed in animal skins. Linnaeus named the orangutan Simia satyrus ("satyr monkey"). He placed the three genera Homo, Simia and Lemur in the order of Primates.
The troglodytes name was used for the chimpanzee by Blumenbach in 1775, but moved to the genus Simia. The orangutan was moved to the genus Pongo in 1799 by Lacépède.
Linnaeus's inclusion of humans in the primates with monkeys and apes was troubling for people who denied a close relationship between humans and the rest of the animal kingdom. Linnaeus's Lutheran archbishop had accused him of "impiety". In a letter to Johann Georg Gmelin dated 25 February 1747, Linnaeus wrote:
Accordingly, Johann Friedrich Blumenbach in the first edition of his Manual of Natural History (1779), proposed that the primates be divided into the Quadrumana (four-handed, i.e. apes and monkeys) and Bimana (two-handed, i.e. humans). This distinction was taken up by other naturalists, most notably Georges Cuvier. Some elevated the distinction to the level of order.
However, the many affinities between humans and other primates – and especially the "great apes" – made it clear that the distinction made no scientific sense. In his 1871 book The Descent of Man, and Selection in Relation to Sex, Charles Darwin wrote:
Changes in taxonomy and terminology
Characteristics
The lesser apes are the gibbon family, Hylobatidae, of sixteen species; all are native to Asia. Their major differentiating characteristic is their long arms, which they use to brachiate through trees. Their wrists are ball and socket joints as an evolutionary adaptation to their arboreal lifestyle. Generally smaller than the African apes, the largest gibbon, the siamang, weighs up to ; in comparison, the smallest "great ape", the bonobo, is .
The superfamily Hominoidea falls within the parvorder Catarrhini, which also includes the Old World monkeys of Africa and Eurasia. Within this grouping, the two families Hylobatidae and Hominidae can be distinguished from Old World monkeys by the number of cusps on their molars; hominoids have five in the "Y-5" molar pattern, whereas Old World monkeys have only four in a bilophodont pattern.
Further, in comparison with Old World monkeys, hominoids are noted for: more mobile shoulder joints and arms due to the dorsal position of the scapula; broader ribcages that are flatter front-to-back; and a shorter, less mobile spine, with greatly reduced caudal (tail) vertebrae—resulting in complete loss of the tail in extant hominoid species. These are anatomical adaptations, first, to vertical hanging and swinging locomotion (brachiation) and, later, to developing balance in a bipedal pose. Note there are primates in other families that also lack tails, and at least one, the pig-tailed langur, is known to walk significant distances bipedally. The front of the ape skull is characterised by its sinuses, fusion of the frontal bone, and by post-orbital constriction.
Distinction from monkeys
Cladistically, apes, catarrhines, and extinct species such as Aegyptopithecus and Parapithecidaea, are monkeys, so one can only specify ape features not present in other monkeys.
Unlike most monkeys, apes do not possess a tail. Monkeys are more likely to be in trees and use their tails for balance. While the great apes are considerably larger than monkeys, gibbons (lesser apes) are smaller than some monkeys. Apes are considered to be more intelligent than monkeys, which are considered to have more primitive brains.
The enzyme urate oxidase has become inactive in all apes, its function having been lost in two primate lineages during the middle Miocene; first in the common ancestors of Hominidae, and later in the common ancestor of Hylobatidae. It has been hypothesized that in both incidents it was a mutation that occurred in apes living in Europe when the climate was getting colder, leading to starvation during winter. The mutation changed the biochemistry of the apes and made it easier to accumulate fat, which allowed the animals to survive longer periods of starvation. When they migrated to Asia and Africa, this genetic trait remained.
Behaviour
Major studies of behaviour in the field were completed on the three better-known "great apes", for example by Jane Goodall, Dian Fossey and Birutė Galdikas. These studies have shown that in their natural environments, the non-human hominoids show sharply varying social structure: gibbons are monogamous, territorial pair-bonders, orangutans are solitary, gorillas live in small troops with a single adult male leader, while chimpanzees live in larger troops with bonobos exhibiting promiscuous sexual behaviour. Their diets also vary; gorillas are foliovores, while the others are all primarily frugivores, although the common chimpanzee hunts for meat. Foraging behaviour is correspondingly variable.
In November 2023, scientists reported, for the first time, evidence that groups of primates, including apes, and, particularly bonobos, are capable of cooperating with each other.
Diet
Apart from humans and gorillas, apes eat a predominantly frugivorous diet, mostly fruit, but supplemented with a variety of other foods. Gorillas are predominantly folivorous, eating mostly stalks, shoots, roots and leaves with some fruit and other foods. Non-human apes usually eat a small amount of raw animal foods such as insects or eggs. In the case of humans, migration and the invention of hunting tools and cooking has led to an even wider variety of foods and diets, with many human diets including large amounts of cooked tubers (roots) or legumes. Other food production and processing methods including animal husbandry and industrial refining and processing have further changed human diets. Humans and other apes occasionally eat other primates. Some of these primates are now close to extinction with habitat loss being the underlying cause.
Cognition
All the non-human hominoids are generally thought of as highly intelligent, and scientific study has broadly confirmed that they perform very well on a wide range of cognitive tests—though there is relatively little data on gibbon cognition. The early studies by Wolfgang Köhler demonstrated exceptional problem-solving abilities in chimpanzees, which Köhler attributed to insight. The use of tools has been repeatedly demonstrated; more recently, the manufacture of tools has been documented, both in the wild and in laboratory tests. Imitation is much more easily demonstrated in "great apes" than in other primate species. Almost all the studies in animal language acquisition have been done with "great apes", and though there is continuing dispute as to whether they demonstrate real language abilities, there is no doubt that they involve significant feats of learning. Chimpanzees in different parts of Africa have developed tools that are used in food acquisition, demonstrating a form of animal culture.
Threats and conservation
All non-human hominoids are rare and threatened with extinction. The eastern hoolock gibbon is the least threatened, only being vulnerable to extinction. Five gibbon species are critically endangered, as are all species of orangutan and gorilla. The remaining species of gibbon, the bonobo, and all four subspecies of chimpanzees are endangered. The chief threat to most of the endangered species is loss of tropical rainforest habitat, though some populations are further imperiled by hunting for bushmeat. The great apes of Africa are also facing threat from the Ebola virus. Currently considered to be the greatest threat to survival of African apes, Ebola infection is responsible for the death of at least one third of all gorillas and chimpanzees since 1990.
All the species of great apes in Africa, are considered endangered. Hunting, logging, agricultural expansion and mining are among the main threats. Recently mining has expanded due to the energy transition. According to researchers "This means that current climate solutions could lead to more industrialization in these places, which could worsen the climate crisis". The Sustainable Critical Minerals Alliance was created for solve problems like this.
| Biology and health sciences | Primates | null |
229917 | https://en.wikipedia.org/wiki/Tupolev%20Tu-154 | Tupolev Tu-154 | The Tupolev Tu-154 (; NATO reporting name: "Careless") is a three-engined, medium-range, narrow-body airliner designed in the mid-1960s and manufactured by Tupolev. A workhorse of Soviet and (subsequently) Russian airlines for several decades, it carried half of all passengers flown by Aeroflot and its subsidiaries (137.5 million/year or 243.8 billion passenger-km in 1990), remaining the standard domestic-route airliner of Russia and former Soviet states until the mid-2000s. It was exported to 17 non-Russian airlines and used as a head-of-state transport by the air forces of several countries.
The aircraft has a cruising speed of and a range of . Capable of operating from unpaved and gravel airfields with only basic facilities, it was widely used in the extreme Arctic conditions of Russia's northern/eastern regions, where other airliners were unable to operate. Originally designed for a 45,000-hour service life (18,000 cycles), but capable of 80,000 hours with upgrades, it was expected to continue in service until 2016, although noise regulations have restricted flights to Western Europe and other regions.
Development
The Tu-154 was developed to meet Aeroflot's requirement to replace the jet-powered Tu-104 and the Antonov An-10 and Ilyushin Il-18 turboprops. The requirements called for either a payload capacity of with a range of while cruising at , or a payload of with a range of while cruising at . A take-off distance of at maximum takeoff weight was also stipulated as a requirement. Conceptually similar to the British Hawker Siddeley Trident, which first flew in 1962, and the American Boeing 727, which first flew in 1963, the medium-range Tu-154 was marketed by Tupolev at the same time as Ilyushin was marketing its long-range Ilyushin Il-62. The Soviet Ministry of Aircraft Industry chose the Tu-154, as it incorporated the latest in Soviet aircraft design and best met Aeroflot's anticipated requirements for the 1970s and 1980s.
The first project chief was ; in 1964, assumed that position. In 1975, the project lead role was turned over to .
The Tu-154 first flew on 4 October 1968. The first deliveries to Aeroflot were in 1970 with freight (mail) services beginning in May 1971 and passenger services in February 1972. Limited production of the 154M model was still occurring as of January 2009, despite previous announcements of the end of production in 2006. In total, 1025 Tu-154s have been built, 214 of which were still in service as of 14 December 2009. The last serial Tu-154 was delivered to the Russian Defense Ministry on 19 February 2013 from the Aviakor factory, equipped with upgraded avionics, a VIP interior, and a communications suite. The factory has four unfinished airframes in its inventory, which can be completed if new orders are received.
Design
The Tu-154 is powered by three rear-mounted, low-bypass turbofan engines arranged similarly to those of the Boeing 727, but it is slightly larger than its American counterpart. Both the 727 and the Tu-154 use an S-duct for the middle (number-two) engine. The original model was equipped with Kuznetsov NK-8-2 engines, which were replaced with Soloviev D-30KU-154s in the Tu-154M. All Tu-154 aircraft models have a relatively high thrust-to-weight ratio, giving the type excellent performance, though at the expense of lower fuel efficiency. This became an important factor in later decades as fuel costs grew.
The cockpit is fitted with conventional dual yoke control columns. Flight control surfaces are hydraulically operated.
The cabin of the Tu-154, although of the same six-abreast seating layout, gives the impression of an oval interior, with a lower ceiling than is common on Boeing and Airbus airliners. The passenger cabin accommodates 128 passengers in a two-class layout and 164 passengers in single-class layout, and up to 180 passengers in high-density layout. The layout can be modified to a winter version where some seats are taken out and a wardrobe is installed for passenger coats. The passenger doors are smaller than on its Boeing and Airbus counterparts. Luggage space in the overhead compartments is very limited.
Like the Tupolev Tu-134, the Tu-154 has a wing swept back at 35° at the quarter-chord line. The British Hawker Siddeley Trident has the same sweepback angle, while the Boeing 727 has a slightly smaller sweepback angle of 32°. The wing also has anhedral (downward sweep) which is a distinguishing feature of Russian low-wing airliners designed during this era. Most Western low-wing airliners such as the contemporary Boeing 727 have dihedral (upward sweep). The anhedral means that Russian airliners have poor lateral stability compared to their Western counterparts, but also are more resistant to Dutch roll tendencies.
Considerably heavier than its predecessor Soviet-built airliner, the Ilyushin Il-18, the Tu-154 was equipped with an oversized landing gear to reduce ground load, enabling it to operate from the same runways. The aircraft has two six-wheel main bogies fitted with large, low-pressure tires that retract into pods extending from the trailing edges of the wings (a common Tupolev feature), plus a two-wheel nose gear unit. Soft oleo struts (shock absorbers) provide a much smoother ride on bumpy airfields than most airliners, which very rarely operate on such poor surfaces.
The original requirement was to have a three-person flight crew – captain, first officer, and flight engineer – as opposed to a four- or five-person crew, as on other Soviet airliners. A fourth crew member, a navigator, was soon found to be still needed, and a seat was added on production aircraft, although that workstation was compromised due to the limitations of the original design. Navigators are no longer trained, and this profession is becoming obsolete with the retirement of the oldest Soviet-era planes.
The latest variant (Tu-154M-100, introduced 1998) includes an NVU-B3 Doppler navigation system, a triple autopilot, which provides an automatic ILS approach according to ICAO category II weather minima, an autothrottle, a Doppler drift and speed measure system, and a "Kurs-MP" radio navigation suite. A stability and control augmentation system improves handling characteristics during manual flight. Modern upgrades normally include modernised TCAS, GPS, and other systems (mostly American- or EU-made).
Early versions of the Tu-154 cannot be modified to meet the current Stage III noise regulations, so are no longer allowed to fly into airspace where such regulations are enforced, such as the European Union, but the Tu-154M's D-30 engines can be fitted with hush kits, allowing them to meet noise regulations.
Variants
Many variants of this airliner have been built. Like its Western counterpart, the Boeing 727, many of the Tu-154s in service have been hush-kitted, and some converted to freighters.
Tu-154
Tu-154 production started in 1970, and the first passenger flight was performed on 9 February 1972. Powered by Kuznetsov NK-8-2 turbofans, it carried 164 passengers. About 42 were built.
Tu-154A
The first upgraded version of the original Tu-154, the A model, in production since 1974, added center-section fuel tanks and more emergency exits, while engines were upgraded to higher-thrust Kuznetsov NK-8-2U. Other upgrades include automatic flaps/slats and stabilizer controls and modified avionics. Max. takeoff weight – . There were 15 different interior layouts for the different domestic and international customers, seating between 144 and 152 passengers. To discern the A model from the base model note the spike at the junction of the fin and tail. This is a fat bullet on the A model, and a slender spike on the base model.
Tu-154B
As the original Tu-154 and Tu-154A suffered wing cracks after a few years in service, a version with a new, stronger wing, designated Tu-154B, went into production in 1975. It also had an additional fuselage fuel tank, additional emergency exits in the tail. Also, the maximum takeoff weight increased to . Important to Aeroflot was the increased passenger capacity, hence lower operating costs. With the NK-8-2U engines the only way to improve the economics of the airplane was to spread costs across more seats. The autopilot was certified for ICAO Category II automatic approaches. Most previously built Tu-154 and Tu-154A were also modified into this variant, with the replacement of the wing. Maximum takeoff weight increased to . 111 were built.
Tu-154B-1
Aeroflot wanted this version for increased revenue on domestic routes. It carried 160 passengers. This version also had some minor modifications to the fuel system, avionics, air conditioning, and landing gear. 64 were built from 1977 to 1978.
Tu-154B-2
A minor modernization of Tu-154B-1. The airplane was designed to be converted from the 160 passenger version to a 180 passenger version by removing the galley. The procedure took about hours. Some of the earlier Tu-154Bs were modified to that standard. Maximum takeoff weight increased to , later to . Some 311 aircraft were built, including VIP versions. A few remain in service.
Tu-154S
The Tu-154S is an all-cargo or freighter version of the Tu-154B, using a strengthened floor, and adding a forward cargo door on the port side of the fuselage. The aircraft could carry nine Soviet PAV-3 pallets. Maximum payload – . There were plans for 20 aircraft, but only nine were converted, two from Tu-154 models and seven from Tu-154B models. Trials were held in the early 1980s and the aircraft was authorized regular operations in 1984. By 1997 all had been retired.
Tu-154M
The Tu-154M and Tu-154M Lux are the most highly upgraded versions, which first flew in 1982 and entered mass production in 1984. It uses more fuel-efficient Soloviev D-30KU-154 turbofans. Together with significant aerodynamic refinement, this led to much lower fuel consumption hence longer range, as well as lower operating costs. The aircraft has new double-slotted (instead of triple-slotted) flaps, with an extra 36-degree position (in addition to existing 15, 28 and 45-degree positions on older versions), which allows reduction of noise on approach. It also has a relocated auxiliary power unit and numerous other improvements. Maximum takeoff weight increased first to , then to . Some aircraft are certified to . About 320 were manufactured. Mass production ended in 2006, though limited manufacturing continued as of January 2009. No new airframes have been built since the early 1990s, and production since then involved assembling aircraft from components on hand. Chinese Tu-154MD electronic intelligence aircraft carry a large-size synthetic-aperture radar (SAR) under their mainframe.
Tu-154M-LK-1
Cosmonaut trainer. This was a salon VIP aircraft modified to train cosmonauts to fly the Buran reusable spacecraft, the Soviet equivalent of the US Space Shuttle. The Tu-154 was used because the Buran required a steep descent, which the Tu-154 was capable of replicating. The cabin featured trainee workstations, one of which was identical to the Buran's flightdeck. The forward baggage compartment was converted into a camera bay, as the aircraft was used to train cosmonauts in observation and photographic techniques.
Tu-154M-ON monitoring aircraft
Germany modified one of the Tu-154s it inherited from the former East German Air Force into an observation airplane. This aircraft was involved with the Open Skies inspection flights. It was converted at the Elbe Aircraft Plant (Elbe Flugzeugwerke) in Dresden, and flew in 1996. After 24 monitoring missions, it was lost in a mid-air collision in 1997.
The Russians also converted a Tu-154M to serve as an Open Skies monitoring aircraft. They used the Tu-154M-LK-1, and converted it to a Tu-154M-ON. When not flying over North America, it is used to ferry cosmonauts. China is believed to have converted one Tu-154 to an electronic countermeasures aircraft.
Tu-154M-100
Design of this variant started in 1994, but the first aircraft were not delivered until 1998. It is an upgraded version with Western avionics, including the Flight Management Computer, GPS, EGPWS, TCAS, and other modern systems. The airplane could carry up to 157 passengers. The cabin featured an automatic oxygen system and larger overhead bins. Three were produced, as payment of debts owed by Russia to Slovakia. Three aircraft were delivered in 1998 to Slovak Airlines, and sold back to Russia in 2003.
Tu-155
A Tu-154 converted into a testbed for alternative fuels. It first flew in 1988 and was used until the fall of the Soviet Union, when it was placed in storage.
Tu-156
Proposed conversions of three Tu-154s with Kuznetsov NK-89 turbofans running on liquid natural gas. Not proceeded with.
Tu-164
Initial designation of the Tu-154M.
Tu-174
Proposed stretched version of Tu-154.
Tu-194
Proposed shortened version of Tu-154.
Operators
Current operators
As of August 2017, there were 44 Tupolev Tu-154 aircraft of all variants still in civil, governmental or military service.
A 45th aircraft has been sighted flying with Air Kyrgyzstan in 2017, but is not listed by the airline as part of its fleet. A 46th aircraft, a Polish Tu-154 with operational number 102, is currently in storage at the military airport in Mińsk Mazowiecki. It was operated by 36th Special Aviation Regiment, but after the 2010 Polish Air Force Tu-154 crash of the Tu-154 101, the Regiment has been disbanded and the plane was grounded. It was fully operational, but the government decided not to use or sell it until the investigation into the Smoleńsk crash is finished. As of June 2021 the aircraft is not flying, and it is unlikely to come back into service, since the government operates a fleet of brand-new, more fuel-efficient jets like the Gulfstream G550 and the Boeing 737 NG. In 2020 it was revealed by the investigation team, led by Antoni Macierewicz, that the aircraft was structurally damaged. The access to the aircraft was restricted by the general prosecutor, and entering its hangar requires a special permission.
As of June 2015, the remaining operators were:
Operational history
In January 2010 Russian flag carrier Aeroflot announced the retirement of its Tu-154 fleet after 40 years, with the last scheduled flight being Aeroflot Flight 736 from Yekaterinburg to Moscow on 31 December 2009. In December 2010, Uzbekistan Airways also declared that it was retiring its Tu-154s. In February 2011, all remaining Iranian Tu-154s were grounded after two incidents.
On 27 December 2016, the Russian Ministry of Defence announced that it had grounded all of its Tu-154s until the end of the investigation into the December 2016 crash of a 1983 Tupolev Tu-154. This was followed by the grounding of all Tu-154s in Russia. The Tu-154 had crashed into the Black Sea just after takeoff from Sochi, Russia, on 25 December 2016 killing all 92 people on board, including 64 members of the Alexandrov Ensemble, an official army choir of the Russian Armed Forces.
In October 2020 ALROSA, the last Russian passenger airline to operate this aircraft, retired its last remaining Tu-154.
Former operators
Former civil operators
Former military operators
Incidents and accidents
Between 1970 and December 2016 there were 110 serious incidents involving the Tu-154, including 73 hull losses, with 2,911 fatalities. Of the fatal incidents, five resulted from terrorist or military terrorist action (two other wartime losses were non-fatal), several from poor runway conditions in winter (including one in which the airplane struck snow plows on the runway), cargo overloading in the lapse of post-Soviet federal safety standards, and mid-air collisions due to faulty air traffic control. Other incidents resulted from mechanical problems, running out of fuel on unscheduled routes, pilot errors (including inadequate flight training for new crews), and cargo fires; several accidents remain unexplained.
On 2 January 2011, Russia's Federal Transport Oversight Agency advised airlines to stop using remaining examples of the Tu-154 (B variant) until the fatal fire incident in Surgut had been investigated. Its operation in Iran ceased in February 2011 due to a number of crashes and incidents involving the type (almost 9% of all Tu-154 losses have occurred in Iran). This grounding compounded the effects of US embargo on civil aircraft parts, substantially decreasing the number of airworthy aircraft in the Iranian civil fleet. In 2010 there were two fatal losses of the Tu-154 due to pilot error and/or weather conditions (a Polish presidential jet attempting a rural airfield landing in heavy fog, the 2010 Polish Air Force Tu-154 crash, and a Russian-registered plane that suffered engine stall after a crew member accidentally de-activated a fuel transfer pump). Following these accidents, in March 2011 the Russian Federal Bureau of Aviation recommended a withdrawal of remaining Tu-154Ms from service.
On 27 December 2016, the Russian Defence Ministry grounded all Tu-154s in Russia pending investigation into the 25 December 2016 Tupolev Tu-154 crash which killed 64 members of the Alexandrov Ensemble, an official Red Army Choir of the Russian Armed Forces.
Aircraft on display
СССР-85020 at the Ukraine State Aviation Museum.
HA-LCG at the Aeropark museum in Budapest.
LZ-BTU at the Aviomuseum Burgas museum in Burgas, Bulgaria.
OK-BYZ at the Aviation Museum Kunovice, Czech Republic
OM-BYO at the Museum of Aviation in Košice, Slovakia
Specifications
In popular culture
Air Crew is the 1979 action film revolving around the exploits of a Soviet Tu-154 crew on an international flight, the first Soviet film in the disaster genre.
| Technology | Specific aircraft_2 | null |
229925 | https://en.wikipedia.org/wiki/Taxus | Taxus | Taxus is a genus of coniferous trees or shrubs known as yews in the family Taxaceae. Yews occur around the globe in temperate zones of the northern hemisphere, northernmost in Norway and southernmost in the South Celebes. Some populations exist in tropical highlands.
The oldest known fossil species are from the Early Cretaceous.
Morphology
They are relatively slow-growing and can be very long-lived, and reach heights of , with trunk girth averaging . They have reddish bark, lanceolate, flat, dark-green leaves long and broad, arranged spirally on the stem, but with the leaf bases twisted to align the leaves in two flat rows either side of the stem.
The male cones are globose, across, and shed their pollen in early spring. Yews are mostly dioecious, but occasional individuals can be variably monoecious, or change sex with time.
The seed cones are highly modified, each cone containing a single seed long partly surrounded by a modified scale which develops into a soft, bright red berry-like structure called an aril, long and wide and open at the end. The arils are mature 6–9 months after pollination, and with the seed contained are eaten by thrushes, waxwings and other birds, which disperse the hard seeds undamaged in their droppings; maturation of the arils is spread over 2–3 months, increasing the chances of successful seed dispersal.
Taxonomy and systematics
Taxus is the Latin word for this tree and its wood that is used to make javelins. The Latin word is probably borrowed, via Greek tóxon, from taxša, the Scythian word used for "yew" and "bow" (cognate of Persian Taxš meaning bow) because the Scythians used its wood to make their bows.
All of the yews are very closely related to each other, and some botanists treat them all as subspecies or varieties of just one widespread species; under this treatment, the species name used is Taxus baccata, the first yew described scientifically.
Taxus species appear similar. Attempts at taxonomy vary from describing all yews as subspecies of T. baccata, as did RKF Pilger in 1903, to splitting species by even very small morphological differences, as did R. W. Spjut in 2007 with 25 species and over 50 varieties. Some species have traditionally been recognized by geographic distribution, but Asian species have been more difficult to classify. Taxus contorta in the Western Himalaya and Taxus sumatrana in Malesia are now generally agreed upon, but overlapping ranges in the Eastern Himalaya, China, and subtropical southeast Asia have led to greater confusion, with the species Taxus chinensis, Taxus mairei, and Taxus wallichiana being elucidated only in the 21st century with the aid of molecular phylogenetics.
The most distinct is the Sumatran yew (T. sumatrana, native to Sumatra and Celebes north to southernmost China), distinguished by its sparse, sickle-shaped yellow-green leaves. The Mexican yew (Taxus globosa, native to eastern Mexico south to Honduras) is also relatively distinct with foliage intermediate between Sumatran yew and the other species. The Florida yew, Mexican yew and Pacific yew are all rare species listed as threatened or endangered.
Distribution
Yews typically occur in the understory or canopy of moist temperate or tropical mountain forests. Elevation varies by latitude from in tropical forests to near sea level in its northernmost populations.
Yews are common in landscape architecture, giving rise to widespread naturalized populations in the United States. There, both T. baccata and Taxus cuspidata are common ornamental shrubs.
T. baccata appears throughout Europe and into western Asia. T. cuspidata occurs over much of East Asia, in China, Japan, Korea, and Sakhalin. Taxus brevifolia ranges in the United States from California to Montana and Alaska, while Taxus canadensis appears in the northeastern United States and southeast Canada.
Species and hybrids
Plants of the World Online recognizes 12 confirmed species:
Taxus baccata , European yew
Taxus brevifolia , Pacific yew, western yew
Taxus calcicola , Asian limestone yew
Taxus canadensis , Canada yew
Taxus chinensis , China yew
Taxus contorta , West Himalayan yew
Taxus cuspidata , Rigid branch yew, Japanese yew
Taxus floridana , Florida yew
Taxus florinii , Florin yew
Taxus globosa , Mesoamerican yew
Taxus mairei , Maire yew
Taxus wallichiana , Wallich yew, East Himalayan yew
Fossil (extinct) species
Taxus engelhardtii – Oligocene, Bohemia, twig-leaves, similar to T. mairei
Taxus inopinata – Upper Miocene, leaf, similar to T. baccata
Taxus masonii – Eocene Clarno Formation; Oregon, USA
Taxus schornii – Miocene, northern Idaho
Commonly reported hybrids
Taxus × media = Taxus baccata × Taxus cuspidata
Taxus × hunnewelliana = Taxus cuspidata × Taxus canadensis
Phylogeny
Below are cladograms showing the evolutionary relationships between yew species and their global distribution.
Toxicity
All species of yew contain highly poisonous taxine alkaloids, with some variation in the exact formula of the alkaloid between the species. All parts of the tree except the arils contain the alkaloid. The arils are edible and sweet, but the seed is dangerously poisonous; unlike birds, the human stomach can break down the seed coat and release the toxins into the body. This can have fatal results if yew 'berries' are eaten without removing the seeds first. Grazing animals, particularly cattle and horses, are also sometimes found dead near yew trees after eating the leaves, though deer are able to break down the poisons and will eat yew foliage freely. In the wild, deer browsing of yews is often so extensive that wild yew trees are commonly restricted to cliffs and other steep slopes inaccessible to deer. The foliage is also eaten by the larvae of some Lepidopteran insects including the moth willow beauty.
Allergenic potential
All parts of a yew plant are toxic to humans with the exception of the yew berries (which however contain a toxic seed); additionally, male and dioecious yews in this genus release cytotoxic pollen, which can cause headaches, lethargy, aching joints, itching, and skin rashes; it is also a trigger for asthma. These pollen granules are extremely small, and can easily pass through window screens. Male yews bloom and release abundant amounts of pollen in the spring; completely female yews only trap pollen while producing none.
Yews in this genus are primarily separate-sexed, and males are extremely allergenic, with an OPALS allergy scale rating of 10 out of 10. Completely female yews have an OPALS rating of 1, and are considered "allergy-fighting".
Uses and traditions
Bows
Yew wood is reddish brown (with whiter sapwood), and is very springy. It was traditionally used to make bows, especially the longbow. These longbows were used by Scythian people who were part of the police force in ancient Athens. This use was lent into the Ancient Greek word for "bow" and later probably borrowed into the Latin word and now generic name of Taxus.
Ötzi, the Chalcolithic mummy found in 1991 in the Italian Alps, carried an unfinished bow made of yew wood. Consequently, it is not surprising that in Norse mythology, the abode of the god of the bow, Ullr, had the name Ydalir (Yew Dales). Most longbow wood used in northern Europe was imported from Iberia, where climatic conditions are better for growing the knot-free yew wood required. The yew longbow was the critical weapon used by the English in the defeat of the French cavalry at the Battle of Agincourt, 1415. British yews tend to be too gnarly, and thus the wood for English longbows used at the Battle of Agincourt was imported from Spain or northern Italy.
Cultivation
It is suggested that English parishes were required to grow yews and, because of the trees' toxic properties, they were grown in the only commonly enclosed area of a village – the churchyard. The yew tree can often be found in church graveyards and is symbolic of sadness. Such a representation appears in Lord Alfred Tennyson's poem "In Memoriam A.H.H." (2.61–64).
The yew can be very long-lived. The Fortingall Yew has been considered to be the oldest tree in Europe, at something over 2,000 years old. Tradition has it that Pontius Pilate slept under it while on duty before 30 AD. Claims for an older tree have been made for the Defynnog Yew in the churchyard of St Cynog's Church, Defynnog, Wales, but this view is contested. Such old trees usually consist of a circular ring of growths of yew, since their heart has long since rotted away.
The Eihwaz rune is named after the yew, and sometimes also associated with the "evergreen" world tree, Yggdrasil.
Horticulture
Yews are widely used in landscaping and ornamental horticulture. Over 400 cultivars of yews have been named, the vast majority of these being derived from European yew (Taxus baccata) or Japanese yew (Taxus cuspidata). The hybrid between these two species is Taxus × media. A popular fastigiate selection of the European yew (Taxus baccata 'Fastigiata') is often called the Irish yew, illustrating the difficulties with common names. A few cultivars with yellow leaves are collectively known as golden yews.
Chemistry
The Pacific yew (Taxus brevifolia), native to the Pacific Northwest of North America, and the Canada yew (Taxus canadensis) of Eastern and Central North America were the initial sources of paclitaxel or Taxol, a chemotherapeutic drug used in breast and lung cancer treatment and, more recently, in the production of the Taxus drug eluting stent by Boston Scientific. Over-harvesting of the Pacific yew for paclitaxel led to fears that it would become an endangered species, since the drug was initially extracted from the bark of the yew, the harvesting of which kills the tree. On January 18, 2008, the Botanic Gardens Conservation International (representing botanic gardens in 120 countries) stated that "400 medicinal plants are at risk of extinction, from over-collection and deforestation, threatening the discovery of future cures for disease." These included yew trees, whose bark is used for the cancer drug paclitaxel.
However, methods were developed to produce the drug semi-synthetically from the leaves of cultivated European yews. Those can be sustainably harvested without the need to further endanger wild populations, and the Pacific yew is no longer at risk. The more common Canada yew is also being successfully harvested in northern Ontario, Quebec and New Brunswick, and has become another major source of paclitaxel. Other yew species contain similar compounds with similar biochemical activity. Docetaxel, an analogue of paclitaxel, is derived from the European yew (Taxus baccata).
In culture
The yew tree is a frequent symbol in the Christian poetry of T. S. Eliot, especially his Four Quartets.
| Biology and health sciences | Pinophyta (Conifers) | Plants |
229940 | https://en.wikipedia.org/wiki/Multiplicative%20inverse | Multiplicative inverse | In mathematics, a multiplicative inverse or reciprocal for a number x, denoted by 1/x or x−1, is a number which when multiplied by x yields the multiplicative identity, 1. The multiplicative inverse of a fraction a/b is b/a. For the multiplicative inverse of a real number, divide 1 by the number. For example, the reciprocal of 5 is one fifth (1/5 or 0.2), and the reciprocal of 0.25 is 1 divided by 0.25, or 4. The reciprocal function, the function f(x) that maps x to 1/x, is one of the simplest examples of a function which is its own inverse (an involution).
Multiplying by a number is the same as dividing by its reciprocal and vice versa. For example, multiplication by 4/5 (or 0.8) will give the same result as division by 5/4 (or 1.25). Therefore, multiplication by a number followed by multiplication by its reciprocal yields the original number (since the product of the number and its reciprocal is 1).
The term reciprocal was in common use at least as far back as the third edition of Encyclopædia Britannica (1797) to describe two numbers whose product is 1; geometrical quantities in inverse proportion are described as in a 1570 translation of Euclid's Elements.
In the phrase multiplicative inverse, the qualifier multiplicative is often omitted and then tacitly understood (in contrast to the additive inverse). Multiplicative inverses can be defined over many mathematical domains as well as numbers. In these cases it can happen that ; then "inverse" typically implies that an element is both a left and right inverse.
The notation f −1 is sometimes also used for the inverse function of the function f, which is for most functions not equal to the multiplicative inverse. For example, the multiplicative inverse is the cosecant of x, and not the inverse sine of x denoted by or . The terminology difference reciprocal versus inverse is not sufficient to make this distinction, since many authors prefer the opposite naming convention, probably for historical reasons (for example in French, the inverse function is preferably called the ).
Examples and counterexamples
In the real numbers, zero does not have a reciprocal (division by zero is undefined) because no real number multiplied by 0 produces 1 (the product of any number with zero is zero). With the exception of zero, reciprocals of every real number are real, reciprocals of every rational number are rational, and reciprocals of every complex number are complex. The property that every element other than zero has a multiplicative inverse is part of the definition of a field, of which these are all examples. On the other hand, no integer other than 1 and −1 has an integer reciprocal, and so the integers are not a field.
In modular arithmetic, the modular multiplicative inverse of a is also defined: it is the number x such that . This multiplicative inverse exists if and only if a and n are coprime. For example, the inverse of 3 modulo 11 is 4 because . The extended Euclidean algorithm may be used to compute it.
The sedenions are an algebra in which every nonzero element has a multiplicative inverse, but which nonetheless has divisors of zero, that is, nonzero elements x, y such that xy = 0.
A square matrix has an inverse if and only if its determinant has an inverse in the coefficient ring. The linear map that has the matrix A−1 with respect to some base is then the inverse function of the map having A as matrix in the same base. Thus, the two distinct notions of the inverse of a function are strongly related in this case, but they still do not coincide, since the multiplicative inverse of Ax would be (Ax)−1, not A−1x.
These two notions of an inverse function do sometimes coincide, for example for the function where is the principal branch of the complex logarithm and :
.
The trigonometric functions are related by the reciprocal identity: the cotangent is the reciprocal of the tangent; the secant is the reciprocal of the cosine; the cosecant is the reciprocal of the sine.
A ring in which every nonzero element has a multiplicative inverse is a division ring; likewise an algebra in which this holds is a division algebra.
Complex numbers
As mentioned above, the reciprocal of every nonzero complex number is complex. It can be found by multiplying both top and bottom of 1/z by its complex conjugate and using the property that , the absolute value of z squared, which is the real number :
The intuition is that
gives us the complex conjugate with a magnitude reduced to a value of , so dividing again by ensures that the magnitude is now equal to the reciprocal of the original magnitude as well, hence:
In particular, if ||z||=1 (z has unit magnitude), then . Consequently, the imaginary units, , have additive inverse equal to multiplicative inverse, and are the only complex numbers with this property. For example, additive and multiplicative inverses of are and , respectively.
For a complex number in polar form , the reciprocal simply takes the reciprocal of the magnitude and the negative of the angle:
Calculus
In real calculus, the derivative of is given by the power rule with the power −1:
The power rule for integrals (Cavalieri's quadrature formula) cannot be used to compute the integral of 1/x, because doing so would result in division by 0:
Instead the integral is given by:
where ln is the natural logarithm. To show this, note that , so if and , we have:
Algorithms
The reciprocal may be computed by hand with the use of long division.
Computing the reciprocal is important in many division algorithms, since the quotient a/b can be computed by first computing 1/b and then multiplying it by a. Noting that has a zero at x = 1/b, Newton's method can find that zero, starting with a guess and iterating using the rule:
This continues until the desired precision is reached. For example, suppose we wish to compute 1/17 ≈ 0.0588 with 3 digits of precision. Taking x0 = 0.1, the following sequence is produced:
x1 = 0.1(2 − 17 × 0.1) = 0.03
x2 = 0.03(2 − 17 × 0.03) = 0.0447
x3 = 0.0447(2 − 17 × 0.0447) ≈ 0.0554
x4 = 0.0554(2 − 17 × 0.0554) ≈ 0.0586
x5 = 0.0586(2 − 17 × 0.0586) ≈ 0.0588
A typical initial guess can be found by rounding b to a nearby power of 2, then using bit shifts to compute its reciprocal.
In constructive mathematics, for a real number x to have a reciprocal, it is not sufficient that x ≠ 0. There must instead be given a rational number r such that 0 < r < |x|. In terms of the approximation algorithm described above, this is needed to prove that the change in y will eventually become arbitrarily small.
This iteration can also be generalized to a wider sort of inverses; for example, matrix inverses.
Reciprocals of irrational numbers
Every real or complex number excluding zero has a reciprocal, and reciprocals of certain irrational numbers can have important special properties. Examples include the reciprocal of e (≈ 0.367879) and the golden ratio's reciprocal (≈ 0.618034). The first reciprocal is special because no other positive number can produce a lower number when put to the power of itself; is the global minimum of . The second number is the only positive number that is equal to its reciprocal plus one:. Its additive inverse is the only negative number that is equal to its reciprocal minus one:.
The function gives an infinite number of irrational numbers that differ with their reciprocal by an integer. For example, is the irrational . Its reciprocal is , exactly less. Such irrational numbers share an evident property: they have the same fractional part as their reciprocal, since these numbers differ by an integer.
The reciprocal function plays an important role in simple continued fractions, which have a number of remarkable properties relating to the representation of (both rational and) irrational numbers.
Further remarks
If the multiplication is associative, an element x with a multiplicative inverse cannot be a zero divisor (x is a zero divisor if some nonzero y, ). To see this, it is sufficient to multiply the equation by the inverse of x (on the left), and then simplify using associativity. In the absence of associativity, the sedenions provide a counterexample.
The converse does not hold: an element which is not a zero divisor is not guaranteed to have a multiplicative inverse.
Within Z, all integers except −1, 0, 1 provide examples; they are not zero divisors nor do they have inverses in Z.
If the ring or algebra is finite, however, then all elements a which are not zero divisors do have a (left and right) inverse. For, first observe that the map must be injective: implies :
Distinct elements map to distinct elements, so the image consists of the same finite number of elements, and the map is necessarily surjective. Specifically, ƒ (namely multiplication by a) must map some element x to 1, , so that x is an inverse for a.
Applications
The expansion of the reciprocal 1/q in any base can also act as a source of pseudo-random numbers, if q is a "suitable" safe prime, a prime of the form 2p + 1 where p is also a prime. A sequence of pseudo-random numbers of length q − 1 will be produced by the expansion.
| Mathematics | Basics | null |
229943 | https://en.wikipedia.org/wiki/Ilyushin%20Il-76 | Ilyushin Il-76 | The Ilyushin Il-76 (; NATO reporting name: Candid) is a multi-purpose, fixed-wing, four-engine turbofan strategic airlifter designed by the Soviet Union's Ilyushin design bureau as a commercial freighter in 1967, to replace the Antonov An-12. It was developed to deliver heavy machinery to remote and poorly served areas. Military versions of the Il-76 have been widely used in Europe, Asia and Africa, including use as an aerial refueling tanker and command center.
The Il-76 has seen extensive service as a commercial freighter for ramp-delivered cargo, especially for outsized or heavy items that cannot be carried by other means. It has also been used as an emergency response transport for civilian evacuations as well as for humanitarian aid and disaster relief around the world. Thanks to its ability to operate from unpaved runways, it has been useful in undeveloped areas. Specialized models have also been produced for aerial firefighting and zero-G training.
Design and development
Origins
The aircraft was conceived by Ilyushin in 1967 to meet a requirement for a freighter able to carry a payload of over a range of in less than six hours, able to operate from short and unprepared airstrips, and capable of coping with the worst weather conditions likely to be experienced in Siberia and the Soviet Union's Arctic regions. It was intended to replace the Antonov An-12. Another project design for a double-decked 250-passenger airliner was cancelled. The Il-76 first flew in .
Production of Il-76s was allocated to the Tashkent Aviation Production Association in Tashkent, Uzbekistan, then a republic of the Soviet Union. Some 860 of the basic transport variants were manufactured. In the 1990s, modernized variants also equipped with Soloviev D-30 turbofan engines were developed (MF, TF), with a cargo compartment long by wide by tall; these larger variants were not produced in significant quantity due to the financial difficulties being experienced by the Russian Air Force, which was the primary operator of the type.
Further development
From 2004 onwards, a number of aircraft in commercial service were modernized to the Il-76TD-90VD version; this involved the adoption of the newly developed PS-90 engine to comply with European noise limitations. In 2005, the People's Republic of China placed an order for 34 new Il-76MDs and four Il-78 tankers. In June 2013, Russian military export agency Rosoboronexport announced an order by China for 12 Il-76MD aircraft.
The Il-76 has also been modified into an airborne refuelling tanker, designated the Il-78, around 50 aircraft having been produced. A variant of the Il-76 also serves as a firefighting waterbomber. Its airframe was used as a base for the Beriev A-50 'Mainstay' AEW&C (airborne early warning and control) aircraft; around 25 aircraft were made. Another application for the type was found in Antarctic support flights and for conducting simulated weightlessness training for cosmonauts (akin to the "Vomit Comet" used by NASA). Beriev and NPO Almaz also developed an airborne laser flying laboratory designated A-60, of which two were built, much of this project's details remaining classified.
Il-76MD-90A
It was announced in 2010 that the production of a modernized Il-76, the Il-76MD-90A (also known as project Il-476 during the design stage), would begin; a proposed new production line would be located in Aviastar's facility in Ulyanovsk, Russia, and be operated in cooperation with the Tashkent works. At that point, the construction of two Il-76MD-90A prototypes had begun at the Ulyanovsk facility. The first Il-76MD-90A was rolled out at Aviastar's Ulyanovsk plant on 16 June 2014. On 29 April 2015, it was reported that the Russian Aerospace Forces received the first Il-76MD-90A built at the Ulyanovsk plant "Aviastar-SP" from the 2012 contract for 39 aircraft. The Russian Ministry of Defence (MoD) received its first serial production Ilyushin Il-76MD-90A airlifter on 2 April 2019. As of late 2024, 27 aircraft are ordered to be delivered in the period up to 2028 and 26 had been built, six in 2023 and six in 2024.
Operational history
The first aircraft was delivered to the Soviet Air Force in June 1974 and subsequently became the main Soviet strategic transport aircraft. From 1976, it was operated by Aeroflot.
Between 1979 and 1991, Soviet Air Force Il-76s made 14,700 flights into Afghanistan, transporting 786,200 servicemen and 315,800 tons of freight. The Il-76 carried 89% of Soviet troops and 74% of the freight that was airlifted. As Afghan rebels were unable to shoot down high-flying Il-76s, their tactics were to try and damage it on takeoff or landing. Il-76s were often hit by shoulder-launched Stinger and Strela heat-seeking missiles and large-calibre machine gun fire, but because the strong airframes were able to take substantial damage and remain operational, the aircraft had a remarkably low attrition rate during this period of conflict. Building on that experience, the bulk of the Canadian Forces equipment into Afghanistan was flown in using civilian Il-76s. In 2006, the Russian Air Force had about 200 Il-76s. Civilian users in Russia have 108.
In 2004, a Chinese People's Liberation Army Air Force (PLAAF) Il-76 carried out a flight mission in Afghanistan, and later in 2011, PLAAF Il-76s were sent to Libya to evacuate Chinese citizens. The two missions were the reported first steps of PLAAF developing long-range transportation capability.
Syrian Air Force Il-76s, operating as civil Syrianair aircraft, have been reportedly used to ship weapons, money, and other cargo from Russia and Iran to Syria, according to a defected Syrian military pilot. Since the start of the war, in April 2011 (and up to July 2012), around 20 military flights have been conducted to and from Tehran, via Iraqi airspace. Further information exposes that since around 2012, Syrian Il-76s have regularly flown to Moscow's Vnukovo Airport to fetch shipments of Syrian banknotes that have been useful to Bashar al-Assad's government to survive international sanctions.
On 30 January 2017, an Il-76 firebomber of the Russian EMERCOM agency was deployed to Chile to assist firefighters. The assignment took 39 days.
All Il-76 transport aircraft in service with the RF Aerospace Forces were to receive anti-missile systems, and aircraft reconfiguration started in spring 2019.
Variants
Prototypes and developmental variants
Il-76TD-90/Il-76MD-90 Engine upgrades to Perm PS-90s.
Il-76 firebomber Firefighting aircraft to drop exploding capsules filled with fire retardant.
Il-76PSD SAR version of Il-76MF
Il-96 Early development of convertible passenger/cargo aircraft, (project only, designation re-used later)
Il-150 Proposed Beriev A-50 with Perm PS-90 engines.
Beriev A-60 Airborne laser weapon testbed (Il-76 version 1A)
Special purpose/research variants
Il-76LL With reinforced wing (at least 3 aircraft) to be used as test-bed aeroplane for engine prototypes flight testing in Gromov Flight Research Institute.
Izdeliye-176 Prototype Il-76PP.
Izdeliye-576
Izdeliye-676 Telemetry and communications relay aircraft, for use during trial programmes (prototype).
Izdeliye-776 Telemetry and communications relay aircraft, for use during trial programmes (prototype).
Izdeliye-976 ("SKIP", Il-976, or Il-76SK) – (СКИП – Самолетный Контрольно-Измерительный Пункт, Airborne Check-Measure-and-Control Center) Il-76/A-50 based range control and missile tracking platform. Initially built to support Raduga Kh-55 cruise missile tests.
Izdeliye-1076 Special mission aircraft for unknown duties.
Izdeliye-1176 ELINT electronic intelligence aircraft, or Il-76-11
Il-76-Tu160 tailplane transporter One-off temporary conversion to support Tu-160 emergency modification programme.
Il-76K/Il-76MDK/Il-76MDK-II Reduced-gravity aircraft for cosmonaut training (dlya podgotovki kosmonavtov), used by Yuri Gagarin Cosmonaut Training Center.
Il-76LL Engine testbed (ooniversahl'naya letayuschchaya laboratoriya).
Il-76PP ECM aircraft, major problems with ECM equipment on the Izdeliye-176 only.
Il-84 Maritime search and rescue aircraft (alternative designation – Il-76PS – poiskovo-spasahtel'nyy), not produced.
Military variants
Il-76D ('D' for "Desantnyi", Десантный – "paratrooper transport") has a gun turret in the tail for defensive purposes.
Il-76M Military transport version, (modifitseerovannyy – modified).
Il-76MD Improved military transport version, (modifitseerovannyy Dahl'ny – modified, long-range).
Il-76MD Skal'pel-MT Mobile Hospital
Il-76M/Il-76MD Built without military equipment but designated as Ms and MDs (Gordon – 'Falsies')
Il-76MD-90 An Il-76MD with quieter and more economical Aviadvigatel PS-90 high-bypass turbofan engines.
Il-76MF Stretched military version with a longer fuselage, PS-90A-76 engines, maximum takeoff weight of and a lift capability of . First flew in 1995. Two built and delivered to the Royal Jordanian Air Force, later sold to the Egyptian Air Force.
Il-76MD-M Modernized Il-76MD for the Russian Aerospace Forces.
Il-76MD-90A An upgraded version with a new glass cockpit, upgraded avionics, new one-piece carbon-fibre wing, and Aviadvigatel PS-90A-76 engines. It was also known as Il-476 while in development. Designated as Il-76-MD-90AE for the export market.
Il-76T/Il-76TD Built as military aircraft but given civilian designations. (Gordon – 'Falsie')
Ilyushin Il-78/Il-78M/Il-78MD-90A Aerial refuelling tanker.
Il-78 MKI A customized version of the Il-78 developed for the Indian Air Force.
Il-82 Airborne Command Post/communications relay aircraft, (alternative designation – Il-76VKP – 'version65S').
Beriev A-50/Beriev A-50M/Beriev A-50I/Beriev A-50E Airborne Early Warning & Control aircraft. Beriev given control over the program.
Beriev A-100 An AEW&C version of the Il-76MD-90A. Currently in development, with at least two prototypes built.
Civil variants
Il-76MGA Initial commercial freighter (two prototypes and 12 production) equipped with Soloviev D-30 turbofan engines.
Il-76MD to Il-76TD conversions Complete removal of military equipment, identified by crude cover over OBIGGS inlet in Starboard Sponson.
Il-76P/Il-76TP/Il-76TDP/Il-76MDP firefighting aircraft. The Il-76 waterbomber is a VAP-2 1.5-hour install/removal tanking kit conversion. The Il-76 can carry up to of water; 3.5 times the capacity of the C-130 Hercules. Since this kit can be installed on any Il-76, the designation Il-76TP, Il-76TDP are also used when those versions of the Il-76 are converted into waterbombers. The Il-76P was first unveiled in 1990.
Il-76T ('T' for Transport, Транспортный) unarmed civil cargo transport version. NATO code-name "Candid-A". It first flew on November 4, 1978.
Il-76TD The civil equivalent of the Il-76MD, first flew in 1982, equipped with Soloviev D-30 turbofan engines.
Il-76TD-90 An Il-76TD with Aviadvigatel PS-90 engines and a partial glass cockpit.
Il-76TD-90VD An Il-76TD with Aviadvigatel PS-90 engines and a partial glass cockpit. It was developed specially for Volga-Dnepr cargo company, which operates five aircraft as of 2021.
Il-76TD-S Civilian mobile hospital, similar to Il-76MD Skal'pel-MT.
Il-76TF Civil transport stretched version with Aviadvigatel PS-90 engines. It is the civil version of the Il-76MF (none produced).
Foreign variants
Beriev A-50E/I For the Indian Air Force. Hosts Israeli Phalcon radar for AEW&C and Aviadvigatel PS-90 engines.
Il-76MD tanker Iraqi Air Force tanker conversions.
KJ-2000 Domestic Chinese airborne early warning and control conversion of Il-76, developed after A-50I was cancelled and currently in service with the armed forces of China.
CFTE engine testbed The China Flight Test Establishment (CFTE) currently operates a flying testbed converted from a Russian-made Il-76MD jet transport aircraft to serve as a flying testbed for future engine development programmes. The first engine to be tested on the aircraft is the WS-10A "Taihang" turbofan, currently being developed as the powerplant for China's indigenous J-10 and J-11 fighter aircraft. Il-76MD #76456, acquired by the AVIC 1 from Russia in the 1990s, is currently based at CFTE's flight test facility at Yanliang, Shaanxi Province.
Baghdad-1 Iraqi development with a radar mounted in the cargo hold enabling it to serve as AEW&C, used in the Iran–Iraq War.
Baghdad-2 Iraqi development (with French assistance) with fibreglass-reinforced plastic radome over the antenna of the Thomson-CSF Tiger G surveillance radar with a maximum detection range of . One was destroyed on the ground during the 1991 Persian Gulf War; two others were flown to Iran where they remained. At least one went into service with the IRIAF. One aircraft crashed following a midair collision with a HESA Saeqeh fighter, during the annual Iranian military parade in Tehran. It can be distinguished from the Beriev A-50 by having the Il-76 navigator windows in the nose, which the A-50 does not.
Operators
Military and civil operators in 38 countries have operated more than 850 Il-76 in large numbers. While Russia is the largest military operator of the Il-76, followed by Ukraine and India, Belarus' TransAVIAexport Airlines is the largest civilian operator.
Military operators
Algerian Air Force – 11 Il-76MD / Il-76TD and 5 Il-78TD in service.
National Air Force of Angola – 7 Il-76MD and Il-76TD in service.
Armenian Air Force – 1 Il-76M and 1 Il-76TD in service.
Azerbaijani Air Forces and Air Defense Troops – 2 Il-76TD in service.
Air Force and Air Defence Forces of Belarus – 2 Il-76MD in service.
People's Liberation Army Air Force – 26 Il-76MD / Il-76TD, 3 Il-78M, and 4 KJ-2000 in service.
Air Force of the Democratic Republic of the Congo
Egyptian Air Force – 2 Il-76MF in service.
Equatorial Guinea National Guard - 1 Il-76TD in service.
Indian Air Force – 17 Il-76MD, 6 Il-78MKI, and 3 A-50EI in service. 2 A-50EI on order.
Islamic Republic of Iran Air Force – 5 Il-76TD in service.
Islamic Revolutionary Guards Corps Aerospace Force – 3 Il-76MD / Il-76TD in service.
Royal Jordanian Air Force – 1 Il-76TD in service.
Russian Aerospace Forces – Approximately 120 Il-76, 20 Il-78, 15 A-50, and 1 A-100 in service. 20 Il-76MD-90A, 31 Il-78M-90A on order. More in reserve storage.
National Guard of Russia – 9 Il-76M / Il-76MD in service.
Border Service of the Federal Security Service of the Russian Federation - 3 in service.
Rossiya - Special Flight Detachment
Sudanese Air Force – 1 Il-76TD in service.
Syrian Air Force – 3 Il-76M / 1 Il-76TD in service as of 2023
Uzbekistan Air and Air Defence Forces – 3 Il-76MD in service.
Former military operators
Iraqi Air Force
Libyan Air Force
Soviet Air Forces – Largest former operator of the type, with hundreds of aircraft of multiple variants in service. Passed on to successor states.
Ukrainian Air Force – 7 non-operational as of 2024
Yemeni Air Force
Air Force of Zimbabwe
Civil operators
Azal Avia Cargo
Silk Way Airlines
Belcanto Airlines
TransAVIAexport Airlines
Government of Kazakhstan
GST Aero
Air Koryo
Abakan Avia
Aviacon Zitotrans
Ministry of Emergency Situations
Volga-Dnepr Airlines
Syrian Arab Airlines operates four, including three Il-76Ms. Shared with Syrian Air Force.
Turkmenistan Airlines operates four Il-76TDs.
ATI Aircompany operates a number of Il-76 models.
Lviv Airlines operates three Il-76MDs.
Ukraine Air Alliance operates four, including one Il-76MD and three Il-76TDs.
Veteran Airlines
Volare Airlines operates three, including two Il-76MDs and one Il-76TD.
Yuzhmashavia operates two Il-76TDs.
The United Nations Humanitarian Air Service has operated several of the type from the early to mid-1990s to now. Most of them are either ex-Aeroflot or one that the Russian Aerospace Forces has lent to the UN.
Meridian, Inc, operating an Il-78 with registration N20NS.
Gulf Aviation Technology and Services operates a number of Il-76 aircraft on charter or lease.
Phoenix Aviation used to operate 2 Il-76TDs.
Avialeasing operates the Il-76 on a charter and lease basis.
Uzbekistan Airways used to operate 14 Il-76TDs.
Former civil operators
Angola Air Charter operated an Il-76.
Air Highnesses used to own and operate Il-76T (EK-76300) on behalf of Aéro-Service before it crashed.
Dvin Airlines used to operate an Il-76TD.
Yerevan-Avia used to operate two Il-76 (EK86724 and EK86817).
Global Aviation Services
Belavia operated the Il-76 before its closure in 1999.
Gomelavia operated five Il-76TD.
Faso Airways operated a Il-76TD.
Imtrec Aviation operated a Laotian registered Il-76.
Penas Air Cargo
Air Congo operated an Il-76TD.
The Republic of the Congo operates an Il-76.
Trans Air Congo has operated an Il-76T.
Cubana de Aviación used to operate two Il-76MDs.
Ecuatorial Cargo operated one Il-76TD.
Express International Cargo
Sun Way has operated the Il-76TD.
Atlant Hungary has operated the Il-76.
Hungarian Ukrainian Air Cargo has operated the Il-76.
Atlas Air has operated at least eight Il-76TD.
Chabahar Air has operated at least two Il-76TD.
Mahan Air has operated the Il-76.
Payam Air operated two Il-76TD.
Qeshm Air operated two Il-76TD (airline disestablished).
Safiran Airlines
Yas Air operated two Il-76TDs (registered as EP-GOL and EP-GOM).
Jordan International Air Cargo – has operated two Il-76MFs delivered in 2011. and operated for the Royal Jordanian Air Force which were sold to Egypt in 2018 and were delivered in 2019.
Air Kazakhstan operated Il-76 aircraft until its closure in 2004.
Kazakhstan Airlines operated the Il-76TD before its closure in 1997.
Sayakhat Airlines operated the Il-76 previously.
Air Almaty operates an Il-76TD for leased operations.
Botir Avia used to operate three, including one Il-76MD and two Il-76TD.
Kyrgyzstan Airlines operates one Il-76TD.
Imtrec aviation of Cambodia used to operate Laos registered Il-76TD.
Inversija operated four, including three Il-76Ts and one Il-76TD.
Jamahiria Air Transport operated five Il-76M and 12 Il-76TD.
Libyan Air Cargo, the cargo division of Libyan Arab Airlines, operates 21, including one Il-76M and 15 Il-76TD.
Transafrica Airlines
Aerocom operated an Il-76MD as well as an Il-76T until as late as January 2005.
Airline Transport operated a number of Il-76 aircraft, losing three in accidents in 2004 and 2005.
Jet Line International used to operate the Il-76.
Tiramavia
Aeroflot operated large numbers of aircraft, especially during Soviet years, often on behalf of the Soviet military. However, none remain in service with the airline.
Air STAN operated an Il-76TD.
ALAK operated Il-76 aircraft before its closure in 1999.
Aviaenergo operated the aircraft, but none remain in service.
Continental Airways has operated the Il-76 in the past, but does not do so currently.
Dacono Air has operated the Il-76.
Domodedovo Airlines has operated the Il-76, but none is currently in service.
East Line used to operate the Il-76.
Ilavia Airline used to operate six, including two Il-76MDs and four Il-76TDs.
KrasAir operated the Il-76, but none is currently in service.
Krylo Airlines operated two Il-76TDs into 2005.
Moscow Airways operated an Il-76TD in the early 1990s.
Novosibirsk Air Enterprise operated the Il-76, but none is currently in service.
Pulkovo Aviation Enterprise operated the Il-76, but none is currently in service.
Red Wings Airlines used to operate two Il-76TDs.
Spair Airlines
Tesis Aviation Enterprise used to operate nine Il-76TDs.
Tyumen Airlines
Uralinteravia
Air Tomisko operated 3 Il-76TDs. Two were leased from GST Aero which had been before in service of Kosmas Air, and one more was added in May 2006.
Kosmas Air operated two Il-76TDs leased from GST Aero.
Aerolift Sierra Leone used to operate Il-76 aircraft for special charter and cargo lift operations.
Aeroflot was the main civil user of the aircraft during the period of the Soviet Union, although many of its aircraft were operated on behalf of the military.
Jet Air Cargo was one of the first civil operators of the Il-76 in Russia other than Aeroflot.
Air West operated at least six aircraft, although it is unclear how many remain in service.
East West Cargo operated a number of Il-76s.
Juba Air Cargo operated the Il-76.
Badr Airlines operated two Il-76s.
Trans Attico operated two Il-76TDs.
Alfa Airlines
Azza Transport used to operate two Il-76TDs.
Green Flag Airlines
Air Service Ukraine operated the Il-76MD.
Air Ukraine and Air Ukraine Cargo operated the aircraft, although none were in service at the time of bankruptcy.
Azov Avia Airlines operated two Il-76MDs.
BSL Airline operated as many as six Il-78s.
Busol Airlines operated the Il-76 before its closure in 1998.
Khors Aircompany operated two Il-76MDs.
Ukrainian Cargo Airways operated 21, including 19 Il-76MDs.
South Airlines
Yemenia operated two Il-76TDs.
Accidents and incidents
As of July 2024, the Aviation Safety Network has tracked 137 incidents involving Il-76 series aircraft resulting in the 1,158 fatalities. 99 have been written off in crashes and other accidents. Some of the most notable incidents can be found here.
On 23 November 1979, a Soviet Air Forces Il-76, registration CCCP-86714, banked left during an approach to Vitebsk Airport. Control of the aircraft was lost and the aircraft crashed, killing the crew of seven; this was the first loss of an Il-76.
Christmas day crash. On 25 December 1979, a Soviet Air Force Il-76 crashed 36 km from Kabul, killing 48 onboard.
1988 Soviet Air Force Il-76 crash. On 11 December 1988, an Aeroflot Il-76 crashed on approach to Leninakan, Armenia killing 77 of the 78 on board. The aircraft was on an air relief operation following the 1988 Armenian earthquake.
Soviet Air Force Ilyushin Il-76 crash. On 18 October 1989, a Soviet Air Force Il-76 (CCCP-76569) crashed in the Caspian Sea off Sumqayit, Azerbaijan following wing separation caused by an engine fire, killing all 57 in Azerbaijan's deadliest air accident. The cause of the engine fire was traced back to a design flaw.
1990s
On 1 February 1990, a Soviet Air Forces Il-76 registration СССР-86021 crashed 14 minutes after takeoff from Panevėžys Air Base, killing all 8 members of the crew.
On 24 May 1991, a Metro Cargo Il-76TD (LZ-INK, named Lugano), crashed near Kermanshah Airport while attempting a forced landing following fuel exhaustion, killing four of ten crew.
On 8 July 1993, a Russian Air Force Il-76M (RA-86039) crashed near Pskov Airport due to loss of control following an unexplained in-flight fire, killing the 11 crew.
1995 Airstan Ilyushin Il-76 hijacking. On 3 August 1995 Taliban-controlled fighter aircraft intercepted an Airstan Ilyushin Il-76TD transport aircraft, and held its seven crew members for over a year before escaping.
On 19 August 1996, Spair Airlines Flight 3601, an Il-76T, crashed while trying to land at Belgrade Nikola Tesla Airport following total electrical failure due to pilot error, killing all 11 occupants on board. The crew had forgotten to turn on the AC/DC converter following engine startup.
Charkhi Dadri mid-air collision. On 12 November 1996, Kazakhstan Airlines Flight 1907, an Il-76, collided in mid-air with Saudia Flight 763 (a Boeing 747) over Charkhi Dadri, India, killing all 349 aboard both aircraft in the deadliest mid-air collision. The Kazakh crew failed to maintain altitude owing to confusion with ATC.
1996 Abakan Ilyushin Il-76 crash. On 27 November 1996, a Russian Air Force Ilyushin Il-76MD, registration RA-78804, flew into the side of a mountain, minutes after it departed Abakan Airport, and crashed from the airport. All 21 occupants on board died in the accident.
On 13 July 1998, ATI Aircompany Flight 2570, an Il-76MD (UR-76424), crashed in the sea shortly after takeoff from Ras Al Khaimah International Airport, killing the eight crew. The aircraft was overloaded and the pilot failed to respond to GPWS warnings.
On 17 July 1998, Air Sofia Flight 701, an Il-78 (UR-UCI) struck a hill on approach to Asmara International Airport, killing all ten on board. The aircraft was leased from Ukrainian Cargo Airways.
2000s
Armed Forces of the Russian Federation Flight 9064. On 2 December 2001, a military flight from Bratsk Airport to Petropavlovsk-Kamchatsky Airport crashed at Novaya Inya, Russia, following an onboard fire, killing 18 on board.
Rus Flight 9633. On 14 July 2001 a cargo flight from Chkalovsky Airport crashed shortly after takeoff, killing 10 persons on board the aircraft.
2003 Iran Ilyushin Il-76 crash. On 19 February 2003, an Ilyushin Il-76 crashed near Kerman, Iran under unspecified reasons (possibly weather-related). The crash killed 275 people, including hundreds of the Iranian Revolutionary Guard. The accident remains the deadliest involving the Il-76.
2003 Ukrainian Cargo Airways Il-76 accident. On 8 May 2003, the rear loading ramp of an Il-76 leased by the Congolese government unexpectedly opened at 10,000 feet after taking off from the capital Kinshasa. Initial reports stated that over 120 policemen and their families had been sucked out in 45 minutes, but 14 people actually died.
2007 Mogadishu TransAVIAexport Airlines Il-76 crash. On 23 March 2007, an Il-76 of TransAVIAexport Airlines registered EW-78849 was reportedly shot down (the official cause is undetermined) in the outskirts of Mogadishu, Somalia, killing all 11 people on board.
On 30 June 2008, an Ababeel Aviation Il-76 crashed while taking off from Khartoum on a relief flight, killing the 4 crew members, the only people on board the plane.
On 2 July 2008, Click Airways Flight 1002, operated using an Ilyushin Il-76TD from Bagram Air Base to Al-Fujairah-Fujairah International Airport, suffered an Uncontained Engine Failure of its no. 3 engine at FL280. The failed engine parts struck the no. 4 engine resulting in its failure, as well as the fuselage and fuel tanks. The flight crew managed to successfully make an emergency landing at Zahedan, Iran. None of the three crew sustained injuries.
2009 Makhachkala Il-76 collision. On 15 January 2009, two Russian Ministry of Interior Il-76MDs were involved in a ground collision at Makhachkala Airport. One of the aircraft, registration RA-76825, was ready to depart and was positioned at the runway end when the other one, RA-76827, came in to land. The wing of the landing aircraft struck the flight deck of RA-76825 and a fire erupted. There were three fatalities in the departing aircraft, out of seven occupants on board. None of the 31 occupants aboard RA-76827 were hurt. RA-76825 was written off as a consequence of the accident.
On 9 March 2009, an Aerolift Il-76 (S9-SAB) crashed into Lake Victoria just after takeoff from Entebbe Airport, Uganda, killing all 11 people on board. Two of the engines had caught fire on takeoff. The aircraft was chartered by Dynacorp on behalf of AMISOM. The accident was investigated by Uganda's Ministry of Transport, which concluded that all four engines were time-expired and that Aerolift's claim that maintenance had been performed to extend their service lives and the certification of this work could not be substantiated.
2009 Iranian Air Force Ilyushin Il-76 accident. On 22 September 2009, Islamic Republic of Iran Air Force Il-76MD Adnan 2 "5-8208" Simorgh crashed near Varamin killing all seven people on board. The crash was possibly the result of a mid-air collision with a Northrop F-5E Tiger II.
2009 Yakutia Ilyushin Il-76 crash. On 1 November 2009, an Il-76 belonging to the Russian Ministry of the Interior crashed near the city of Mirny within 2 kilometers after taking off. Eleven people on board were confirmed as killed.
2010s
Sun Way Flight 4412. On 28 November 2010, Il-76 4L-GNI crashed in a populated area of Karachi, Pakistan, shortly after taking off from Jinnah International Airport. All eight people on board were killed, along with two people on the ground. The aircraft was reported to have been trying to return to Jinnah after suffering an uncontained engine failure and fire.
Silk Way Airlines Flight 995. On 6 July 2011, an Il-76, tail number 4K-AZ55, crashed into a mountain in Afghanistan, while on final to Bagram Air Force Base. Eight people on board were initially confirmed as killed, with one unaccounted.
2012 Aéro-Service Ilyushin Il-76 crash. On 30 November 2012, an Aéro-Service Il-76T (also reported as being operated by Trans Air Congo in the days after the accident) crashed 850 meters short of runway 5L of the Congo's Maya-Maya Airport in Brazzaville while landing during a violent storm, killing 32, including the 5 aircrew, another person on board and 26 people on the ground.
2014 Ukrainian Air Force Il-76 shootdown. On 14 June 2014 an Ilyushin Il-76 transport aircraft of the 25th Transport Aviation Brigade of the Ukrainian Air Force was shot down on approach to land at Luhansk International Airport, Ukraine.
2016 Russian Ministry of Emergency Situations Il-76 crash. On 1 July 2016, a Russian Ministry of Emergency Situations (EMERCOM) Il-76TD (RA-76840) struck a hillside near Rybnyi Uyan while fighting wildfires near Irkutsk, killing all ten on board.
2018 Algerian Air Force Ilyushin Il-76 crash. On 11 April 2018, Algerian Air Force Ilyushin Il-76 7T-WIV crashed shortly after take-off from Boufarik Airport, Boufarik, Algeria. All 257 people on board were killed, making the accident the deadliest air crash on Algerian soil.
2020s
25 February 2022, during the 2022 Russian invasion of Ukraine, the Ukrainian State Special Communications Agency and US officials claimed that Russian Il-76s were shot down over Bila Tserkva. As of September 2022, no wreckage of the planes has been found.
4 April 2022, photographs of two destroyed Il-76s from the Ukrainian 25th Transport Aviation Brigade were displayed; these cargo planes were destroyed on the ground by Russian forces at Melitopol Airport.
2022 Russian Air Force Ilyushin Il-76 crash. On 24 June 2022, Russian Aerospace Forces Il-76MD RF-78778 crashed and caught fire whilst landing near the city of Ryazan following an engine fire, killing five of nine on board.
30 August 2023, four Il-76s were reportedly destroyed by Ukrainian kamikaze drone strikes at Pskov Airport.
2023 Gao Ilyushin Il-76 crash. On 23 September 2023, an Il-76 operated by the Malian Armed Forces crashed upon landing at Gao Airport, Mali. According to the French newspaper Le Monde, Malian officials confirmed the aircraft's being owned by the Army and having Wagner Group members on board. The aircraft overshot more than half of the available runway before touching down. For reasons unknown the crew failed to execute a Go-Around in due time which led to the aircraft rolling down the embankment at the end of the runway. The aircraft exploded killing all personnel on board.
On the night of 19–20 October 2023, a Il-76MD military transport plane caught fire during take-off from a military airfield in Dushanbe, the capital of Tajikistan. A wheel exploded on the plane during acceleration, causing a fire to break out. The plane rolled out of the runway and burned down completely. It is known that there were eight people on board. The crew was not injured.
2024 Korochansky Ilyushin Il-76 crash. On 24 January 2024, a Russian Air Force Ilyushin Il-76 crashed in the Korochansky District in the Belgorod Oblast of Russia. The Russian Ministry of Defense stated it was carrying 65 Ukrainian Armed Forces POW, with 6 crew members and 3 security forces at the time. This was refuted by Ukrainian sources citing flight direction, photos of the crash site and other classified information leaked from Russia.
2024 Ivanovo Ilyushin Il-76 crash. On 12 March 2024, an Il-76 crashed in Ivanovo oblast. According to RIA Novosti, the engine caught on fire after the take-off from the Ivanovo air base, and the aircraft crashed when attempting an emergency landing back at the air base. There were eight crew members and seven passengers. All eight crew and seven passengers were killed in the crash.
2024 New Way Cargo Airlines Ilyushin Il-76 shootdown. On 21 October 2024, an Il-76 was shot down over Sudan killing all five crew members onboard.
Aircraft on display
CCCP-76511 (c/n 083414444) preserved in the Ukraine State Aviation Museum, Kyiv. The aircraft was originally painted as UR-UCI of Ukrainian Cargo Airways to commemorate the real aircraft that crashed in 1998, but was returned to its original Aeroflot livery as CCCP-76511 in 2016.
Specifications (Il-76TD)
| Technology | Specific aircraft_2 | null |
229962 | https://en.wikipedia.org/wiki/Feral%20pigeon | Feral pigeon | Feral pigeons (Columba livia domestica or Columba livia forma urbana), also called city doves, city pigeons, or street pigeons, are descendants of domestic pigeons (Columba livia domestica) that have returned to the wild. The domestic pigeon was originally bred from the wild rock dove, which naturally inhabits sea-cliffs and mountains. Rock, domestic, and feral pigeons are all the same species and will readily interbreed. Feral pigeons find the ledges of buildings to be a substitute for sea cliffs, have become adapted to urban life, and are abundant in towns and cities throughout much of the world.
Owing to their capacity to create large amounts of excrement and be an occasional disease vector to humans combined with crop and property damage, pigeons are largely considered a nuisance and an invasive species, often disparagingly referred to as "rats with wings". Actions are taken in many municipalities to lower their numbers or completely eradicate them.
Physical characteristics
Feral pigeons are essentially the same size and shape as the original wild rock dove, but often display far greater variation in colour and pattern than their wild ancestors. The blue-barred pattern which the original wild rock dove displays is generally less common in more urban areas. Urban pigeons tend to have darker plumage than those in more rural areas.
Genetics
The avoidance of mating between related individuals is ordinarily regarded as adaptive since it decreases the likelihood of inbreeding depression in progeny that can be caused by the expression of deleterious recessive alleles. However in feral pigeons it was found that despite detectable inbreeding depression, pairwise relatedness between mates was significantly greater than it was between nonmates. This suggests that mating with close kin provides inclusive fitness benefits that outweigh the costs of inbreeding depression.
Protection status
In the United Kingdom, pigeons are covered under the "General Licences" and can be humanely culled by the land owner or their agent for a variety of reasons including spread of human disease. It is illegal to kill/destroy nests for any reason other than those listed under the general licences.
In the U.S., the Migratory Bird Treaty Act of 1918, which protects native birds, does not apply to feral pigeons, common starlings or house sparrows, because they are introduced species. It is usually legal to kill feral pigeons in the United States; methods such as poisons may be regulated, however.
In India, pigeons are protected under Section 428 and Section 429 of the Indian Penal Code. Wild pigeons are further protected under the Wildlife Protection Act, 1972.
Population control
Feral pigeons often only have small populations within cities relative to the number of humans. For example, the breeding population of feral pigeons in Sheffield, England in summer 2005 was estimated at 12,130 individuals (95% confidence interval 7757–18,970), in a city with a human population of about 500,000. Despite this, feral pigeons usually reach their highest densities in the central portions of cities, so they are frequently encountered by people, which may lead to conflict.
Potential health risk to humans
Feral pigeons are widely considered pests, and can be reservoirs and vectors of some human and livestock diseases, such as salmonellosis and tuberculosis. However, it is rare that a pigeon will transmit a disease to humans due to their immune system. Although feral pigeons pose sporadic health risks to humans, the risk is low, even for humans involved in occupations that bring them into close contact with nesting sites. Analysis revealed that feral pigeons harbored a total of 60 different human pathogenic organisms. Five pathogens were viruses, nine were bacteria, 45 were fungi, and one was a protozoan. However, only five pathogens were routinely transmitted to humans. There were single case incidences for transmission of Salmonella enterica.
Property damage
Pigeons often cause significant pollution with their droppings, though there is little evidence of them driving out other bird species. Pigeons are labeled an invasive species in North America by the USDA.
Predators
Peregrine falcons, which are also originally cliff dwellers, have too adapted to the skyscrapers of large cities and often feed exclusively on rock doves. Some cities actively encourage this through falcon breeding programs. Projects include the Unibase Falcon Project and the Victorian Peregrine Project.
Other predators of the pigeon have been recorded, including Eurasian sparrowhawks, crows and gulls. In London, the population of Great white pelicans at St. James's Park have also been recorded killing and consuming pigeons even when alternative food sources are available. In cities in Western Europe, European herring gulls may occasionally hunt and consume feral pigeons in addition to other birds and small mammals.
Larger birds of prey occasionally take advantage of this population as well. In New York City, the abundance of feral pigeons (and other small animals) has created such a conducive environment for predators that the red-tailed hawk has begun to return in very small numbers, including the notable Pale Male.
Poison
Due to their non-selective nature, most avian poisons have been banned. In the United States market, only 4-aminopyridine (Avitrol) and DRC-1339 remain registered by EPA. DRC-1339 is limited to USDA use only, while 4-AP is a restricted-use pesticide, for use only by licensed applicators.
The use of poisons has been proven to be fairly ineffective, however, as pigeons can breed very quickly, and their numbers are determined by how much food is available; that is, they breed more often when more food is provided to them. When pigeons are poisoned, surviving birds do not leave the area. On the contrary, they are left with more food per bird than before. This attracts pigeons from outside areas as well as encouraging more breeding, and populations are re-established quickly. An additional problem with poisoning is that it also kills pigeon predators. Due to this, in cities with peregrine falcon programs it is typically illegal to poison pigeons.
Reducing food supply
A more effective tactic to reduce the number of feral pigeons is deprivation. Cities around the world have discovered that not feeding their local birds results in a steady population decrease in only a few years. As scavengers, pigeons will still pick at garbage bags containing discarded food or at leftovers carelessly dropped on the ground, but securely disposing of foodstuffs will greatly reduce scavenger populations. Feeding of pigeons is banned in parts of Venice, Italy.
Long-term reduction of feral pigeon populations can be achieved by restricting food supply, which in turn involves legislation and litter (garbage) control. Some cities have deliberately established favorable nesting places for pigeons—nesting places that can easily be reached by city workers who regularly remove eggs, thereby limiting their reproductive success. In addition, pigeon populations may be reduced by bird control systems that successfully reduce nesting sites.
Avian contraceptives
In 1998, in response to conservation groups and the public interest, the National Wildlife Research Center (NWRC), a USDA/APHIS laboratory in Fort Collins, Colorado, started work on nicarbazin, a promising compound for avian contraception. Originally developed for use in resident Canada geese, nicarbazin was introduced for use as a contraceptive for feral pigeons in 2007.
The active ingredient, nicarbazin, interferes with the viability of eggs by binding the ZP-3 sperm receptor site in the egg. This unique contraceptive action is non-hormonal and fully reversible.
Registered by the EPA as a pesticide (EPA Reg. No. 80224-1), "OvoControl P", brand of nicarbazin, is increasingly used in urban areas and industrial sites to control pigeon populations. Declared safe and humane, the new technology is environmentally benign and does not represent a secondary toxicity hazard to raptors or scavengers.
Avian contraception has the support of a range of animal welfare groups including the Humane Society of the United States, the American Society for the Prevention of Cruelty to Animals and People for the Ethical Treatment of Animals. Avian contraceptives are also perceived by some civilians as an acceptable method for population control, over other methods such as prohibition to feeding or extermination.
Dummy egg nesting
When eggs are removed in artificial pigeon houses, the interval between reproductive attempts is strongly reduced, which reduces the efficiency of the method. Dummy egg nesting programs have therefore been tested in some cities with mixed results. There, the eggs are removed and replaced with dummy eggs. The real eggs are then destroyed. One such structure, in Batman Park in Melbourne, Australia, was unsuccessful in attracting pigeons and has since been removed. The loft used in Melbourne was on stilts, with a cage door allowing access from beneath for accessing the structure at night when the pigeons are asleep.
Monitoring pigeon population
Estimating the population size of pigeons is necessary for monitoring and control programs of pigeons in parks and other urban areas. The methods used for estimating populations sizes are:
Stratified grids: This method consists in dividing the area where pigeons occur in 500x500m squares. 34% of the squares are selected randomly and pigeons are counted in a 5 meters radius for 5 minutes.
Point-counts: standing in the center of a park, the observer makes a 360 degree turn while counting individuals with a manual mechanical counter in a radius of approximately 50m, limited by the streets and buildings that surround the park.
Panoramas: taking 360 panoramic photographs, while standing at the center of the park, and using software to place a number above the counted pigeon in the panoramic photograph. This method has been proven the most effective of all.
City squares famous for pigeons
Many city squares have large pigeon populations, such as Washington Square Park in New York City, George Square in Glasgow, the Piazza San Marco in Venice, Dam Square in Amsterdam, The Gateway of India and Kabutarkhana in Mumbai and (prior to 2000) Trafalgar Square in London.
A tall statue of a pigeon by artist Paul Sloan was installed at the Rundle Mall, Adelaide, South Australia, adding to their collection of art installations, including statues of pigs. Sloan intended to "elevate the humble pigeon" with his work titled Pigeon. The mirrored stainless steel statue cost AU$174,000. While the installation has been talked up by City of Adelaide Lord Mayor Sandy Verschoor, some locals have responded negatively.
| Biology and health sciences | Columbimorphae | Animals |
229988 | https://en.wikipedia.org/wiki/Wolf%E2%80%93Rayet%20star | Wolf–Rayet star | Wolf–Rayet stars, often abbreviated as WR stars, are a rare heterogeneous set of stars with unusual spectra showing prominent broad emission lines of ionised helium and highly ionised nitrogen or carbon. The spectra indicate very high surface enhancement of heavy elements, depletion of hydrogen, and strong stellar winds. The surface temperatures of known Wolf–Rayet stars range from 20,000 K to around 210,000 K, hotter than almost all other kinds of stars. They were previously called W-type stars referring to their spectral classification.
Classic (or population I) Wolf–Rayet stars are evolved, massive stars that have completely lost their outer hydrogen and are fusing helium or heavier elements in the core. A subset of the population I WR stars show hydrogen lines in their spectra and are known as WNh stars; they are young extremely massive stars still fusing hydrogen at the core, with helium and nitrogen exposed at the surface by strong mixing and radiation-driven mass loss. A separate group of stars with WR spectra are the central stars of planetary nebulae (CSPNe), post-asymptotic giant branch stars that were similar to the Sun while on the main sequence, but have now ceased fusion and shed their atmospheres to reveal a bare carbon-oxygen core.
All Wolf–Rayet stars are highly luminous objects due to their high temperatures—thousands of times the bolometric luminosity of the Sun () for the CSPNe, for the population I WR stars, to for the WNh stars—although not exceptionally bright visually since most of their radiation output is in the ultraviolet.
The naked-eye star systems γ Velorum and θ Muscae both contain Wolf-Rayet stars, and one of the most massive known stars, R136a1 in 30 Doradus, is also a Wolf–Rayet star.
Observation history
In 1867, using the 40 cm Foucault telescope at the Paris Observatory, astronomers Charles Wolf and Georges Rayet
discovered three stars in the constellation Cygnus (HD 191765, HD 192103 and HD 192641, now designated as WR 134, WR 135, and WR 137 respectively) that displayed broad emission bands on an otherwise continuous spectrum.
Most stars only display absorption lines or bands in their spectra, as a result of overlying elements absorbing light energy at specific frequencies, so these were clearly unusual objects.
The nature of the emission bands in the spectra of a Wolf–Rayet star remained a mystery for several decades. E.C. Pickering theorized that the lines were caused by an unusual state of hydrogen, and it was found that this "Pickering series" of lines followed a pattern similar to the Balmer series when half-integer quantum numbers were substituted. It was later shown that these lines resulted from the presence of helium, the chemical element having just been discovered in 1868.
Pickering noted similarities between Wolf–Rayet spectra and nebular spectra, and this similarity led to the conclusion that some or all Wolf–Rayet stars were the central stars of planetary nebulae.
By 1929, the width of the emission bands was being attributed to Doppler broadening, and hence the gas surrounding these stars must be moving with velocities of 300–2400 km/s along the line of sight. The conclusion was that a Wolf–Rayet star is continually ejecting gas into space, producing an expanding envelope of nebulous gas. The force ejecting the gas at the high velocities observed is radiation pressure.
It was well known that many stars with Wolf–Rayet type spectra were the central stars of planetary nebulae, but also that many were not associated with an obvious planetary nebula or any visible nebulosity at all.
In addition to helium, Carlyle Smith Beals identified emission lines of carbon, oxygen and nitrogen in the spectra of Wolf–Rayet stars.
In 1938, the International Astronomical Union classified the spectra of Wolf–Rayet stars into types WN and WC, depending on whether the spectrum was dominated by lines of nitrogen or carbon-oxygen respectively.
In 1969, several CSPNe with strong oxygen VI (OVI) emissions lines were grouped under a new "OVI sequence", or just OVI type. Similar stars not associated with planetary nebulae were described shortly after and the WO classification was adopted for them. The OVI stars were subsequently classified as [WO] stars, consistent with the population I WR stars.
The understanding that certain late, and sometimes not-so-late, WN stars with hydrogen lines in their spectra are at a different stage of evolution from hydrogen-free WR stars has led to the introduction of the term WNh to distinguish these stars generally from other WN stars. They were previously referred to as WNL stars, although there are late-type WN stars without hydrogen as well as WR stars with hydrogen as early as WN5.
Classification
Wolf–Rayet stars were named on the basis of the strong broad emission lines in their spectra, identified with helium, nitrogen, carbon, silicon, and oxygen, but with hydrogen lines usually weak or absent. Initially simply referred to as class W or W-type stars, the classification was then split into stars with dominant lines of ionised nitrogen (NIII, NIV, and NV) and those with dominant lines of ionised carbon (CIII and CIV) and sometimes oxygen (OIII – OVI), referred to as WN and WC respectively.
The two classes WN and WC were further split into temperature sequences WN5–WN8 and WC6–WC8 based on the relative strengths of the 541.1 nm HeII and 587.5 nm HeI lines. Wolf–Rayet emission lines frequently have a broadened absorption wing (P Cygni profile) suggesting circumstellar material. A WO sequence has also been separated from the WC sequence for even hotter stars where emission of ionised oxygen dominates that of ionised carbon, although the actual proportions of those elements in the stars are likely to be comparable. WC and WO spectra are formally distinguished based on the presence or absence of CIII emission. WC spectra also generally lack the OVI lines that are strong in WO spectra.
The WN spectral sequence was expanded to include WN2–WN9, and the definitions refined based on the relative strengths of the NIII lines at 463.4–464.1 nm and 531.4 nm, the NIV lines at 347.9–348.4 nm and 405.8 nm, and the NV lines at 460.3 nm, 461.9 nm, and 493.3–494.4 nm.
These lines are well separated from areas of strong and variable He emission and the line strengths are well correlated with temperature. Stars with spectra intermediate between WN and Ofpe have been classified as WN10 and WN11 although this nomenclature is not universally accepted.
The type WN1 was proposed for stars with neither NIV nor NV lines, to accommodate Brey 1 and Brey 66 which appeared to be intermediate between WN2 and WN2.5.
The relative line strengths and widths for each WN sub-class were later quantified, and the ratio between the 541.1 nm HeII and 587.5 nm, HeI lines was introduced as the primary indicator of the ionisation level and hence of the spectral sub-class. The need for WN1 disappeared and both Brey 1 and Brey 66 are now classified as WN3b. The somewhat obscure WN2.5 and WN4.5 classes were dropped.
The WC spectral sequence was expanded to include WC4–WC11, although some older papers have also used WC1–WC3. The primary emission lines used to distinguish the WC sub-types are CII 426.7 nm, CIII at 569.6 nm, CIII/IV 465.0 nm, CIV at 580.1–581.2 nm, and the OV (and OIII) blend at 557.2–559.8 nm. The sequence was extended to include WC10 and WC11, and the subclass criteria were quantified based primarily on the relative strengths of carbon lines to rely on ionisation factors even if there were abundance variations between carbon and oxygen.
For WO-type stars the main lines used are CIV at 580.1 nm, OIV at 340.0 nm, OV (and OIII) blend at 557.2–559.8 nm, OVI at 381.1–383.4 nm, OVII at 567.0 nm, and OVIII at 606.8 nm. The sequence was expanded to include WO5 and quantified based the relative strengths of the OVI/CIV and OVI/OV lines.
A later scheme, designed for consistency across classical WR stars and CSPNe, returned to the WO1 to WO4 sequence and adjusted the divisions.
Detailed modern studies of Wolf–Rayet stars can identify additional spectral features, indicated by suffixes to the main spectral classification:
h for hydrogen emission;
ha for hydrogen emission and absorption;
o for no hydrogen emission;
w for weak lines;
s for strong lines;
b for broad strong lines;
d for dust (occasionally vd, pd, or ed for variable, periodic, or episodic dust).
The classification of Wolf–Rayet spectra is complicated by the frequent association of the stars with dense nebulosity, dust clouds, or binary companions. A suffix of "+OB" is used to indicate the presence of absorption lines in the spectrum likely to be associated with a more normal companion star, or "+abs" for absorption lines with an unknown origin.
The hotter WR spectral sub-classes are described as early and the cooler ones as late, consistent with other spectral types. WNE and WCE refer to early type spectra while WNL and WCL refer to late type spectra, with the dividing line approximately at sub-class six or seven. There is no such thing as a late WO-type star. There is a strong tendency for WNE stars to be hydrogen-poor while the spectra of WNL stars frequently include hydrogen lines.
Spectral types for the central stars of planetary nebulae are qualified by surrounding them with square brackets (e.g. [WC4]). They are almost all of the WC sequence with the known [WO] stars representing the hot extension of the carbon sequence. There are also a small number of [WN] and [WC/WN] types, only discovered quite recently.
Their formation mechanism is as yet unclear. Temperatures of the planetary nebula central stars tend to the extremes when compared to population I WR stars, so [WC2] and [WC3] are common and the sequence has been extended to [WC12]. The [WC11] and [WC12] types have distinctive spectra with narrow emission lines and no HeII and CIV lines.
Certain supernovae observed before their peak brightness show WR spectra.
This is due to the nature of the supernova at this point: a rapidly expanding helium-rich ejecta similar to an extreme Wolf–Rayet wind. The WR spectral features only last a matter of hours, the high ionisation features fading by maximum to leave only weak neutral hydrogen and helium emission, before being replaced with a traditional supernova spectrum. It has been proposed to label these spectral types with an "X", for example XWN5(h). Similarly, classical novae develop spectra consisting of broad emission bands similar to a Wolf–Rayet star. This is caused by the same physical mechanism: rapid expansion of dense gases around an extremely hot central source.
Slash stars
The separation of Wolf–Rayet stars from spectral class O stars of a similar temperature depends on the existence of strong emission lines of ionised helium, nitrogen, carbon, and oxygen, but there are a number of stars with intermediate or confusing spectral features. For example, high-luminosity O stars can develop helium and nitrogen in their spectra with some emission lines, while some WR stars have hydrogen lines, weak emission, and even absorption components. These stars have been given spectral types such as O3If∗/WN6 and are referred to as slash stars.
Class O supergiants can develop emission lines of helium and nitrogen, or emission components to some absorption lines. These are indicated by spectral peculiarity suffix codes specific to this type of star:
f for N and He emission
f* for N and He emission with N stronger than N
f+ for emission in Si in addition to N and He
parentheses indicating He absorption lines instead of emission, e.g. (f)
double parentheses indicating strong He absorption and N emission diluted, e.g. ((f+))
These codes may also be combined with more general spectral type qualifiers such as p or a. Common combinations include OIafpe and OIf*, and Ofpe. In the 1970s, it was recognised that there was a continuum of spectra from pure absorption class O to unambiguous WR types, and it was unclear whether some intermediate stars should be given a spectral type such as O8Iafpe or WN8-a. The slash notation was proposed to deal with these situations, and the star Sk−67°22 was assigned the spectral type O3If*/WN6-A. The criteria for distinguishing OIf*, OIf*/WN, and WN stars have been refined for consistency. Slash star classifications are used when the Hβ line has a P Cygni profile; this is an absorption line in O supergiants and an emission line in WN stars. Criteria for the following slash star spectral types are given, using the nitrogen emission lines at 463.4–464.1 nm, 405.8 nm, and 460.3–462.0 nm, together with a standard star for each type:
Another set of slash star spectral types is in use for Ofpe/WN stars. These stars have O supergiant spectra plus nitrogen and helium emission, and P Cygni profiles. Alternatively they can be considered to be WN stars with unusually low ionisation levels and hydrogen. The slash notation for these stars was controversial and an alternative was to extend the WR nitrogen sequence to WN10 and WN11 Other authors preferred to use the WNha notation, for example WN9ha for WR 108. A recent recommendation is to use an O spectral type such as O8Iaf if the 447.1 nm He line is in absorption and a WR class of WN9h or WN9ha if the line has a P Cygni profile. However, the Ofpe/WN slash notation as well as WN10 and WN11 classifications continue to be widely used.
A third group of stars with spectra containing features of both O class stars and WR stars has been identified. Nine stars in the Large Magellanic Cloud have spectra that contain both WN3 and O3V features, but do not appear to be binaries. Many of the WR stars in the Small Magellanic Cloud also have very early WN spectra plus high excitation absorption features. It has been suggested that these could be a missing link leading to classical WN stars or the result of tidal stripping by a low-mass companion.
Nomenclature
The first three Wolf–Rayet stars to be identified, coincidentally all with hot O-class companions, had already been numbered in the Henry Draper catalogue. These stars and others were referred to as Wolf–Rayet stars from their initial discovery but specific naming conventions for them would not be created until 1962 in the "fourth" catalogue of galactic Wolf–Rayet stars. The first three catalogues were not specifically lists of Wolf–Rayet stars and they used only existing nomenclature.
The fourth catalogue of Wolf-Rayet stars numbered them sequentially in order of right ascension. The fifth catalogue used the same numbers prefixed with MR after the author of the fourth catalogue, plus an additional sequence of numbers prefixed with LS for new discoveries. Neither of these numbering schemes remains in common use.
Modern WR catalogues
The sixth Catalogue of Galactic Wolf–Rayet stars was the first to actually bear that name, as well as to describe the previous five catalogues by that name. It also introduced the WR numbers widely used ever since for galactic WR stars. These are again a numerical sequence from WR 1 to WR 158 in order of right ascension.
Compiled in 2001, the seventh catalogue and its annex used the same numbering scheme and inserted new stars into the sequence using lower case letter suffixes, for example WR 102ka for one of the numerous WR stars discovered in the galactic centre. Modern high volume identification surveys use their own numbering schemes for the large numbers of new discoveries. A 2006 Annex was added to the seventh catalog.
In 2011, an online Galactic Wolf Rayet Catalogue was set up, hosted by the University of Sheffield. As of January 2025, it includes 679 stars.
Other numbering schemes
Wolf–Rayet stars in external galaxies are numbered using different schemes. In the Large Magellanic Cloud, the most widespread and complete nomenclature for WR stars is from "The Fourth Catalogue of Population I Wolf–Rayet stars in the Large Magellanic Cloud" prefixed by , for example . Many of these stars are also referred to by their third catalogue number, for example Brey 77. As of 2018, 154 WR stars are catalogued in the LMC, mostly WN but including about twenty-three WCs as well as three of the extremely rare WO class. Many of these stars are often referred to by their RMC (Radcliffe observatory Magellanic Cloud) numbers, frequently abbreviated to just R, for example R136a1.
In the Small Magellanic Cloud SMC WR numbers are used, usually referred to as AB numbers, for example AB7. There are only twelve known WR stars in the SMC, a very low number thought to be due to the low metallicity of that galaxy
In 2012, an IAU working group expanded the numbering system from the Catalogue of Galactic Wolf–Rayet stars so that additional discoveries are given the closest existing WR number plus a numeric suffix in order of discovery. This applies to all discoveries since the 2006 annex, although some of these have already been named under the previous nomenclature; thus WR 42e is now numbered WR 42-1.
Properties
Wolf–Rayet stars are a normal stage in the evolution of very massive stars, in which strong, broad emission lines of helium and nitrogen ("WN" sequence), carbon ("WC" sequence), and oxygen ("WO" sequence) are visible. Due to their strong emission lines they can be identified in nearby galaxies. About 600 Wolf–Rayets have been catalogued in our own Milky Way Galaxy. This number has changed dramatically during the last few years as the result of photometric and spectroscopic surveys in the near-infrared dedicated to discovering this kind of object in the Galactic plane. It is expected that there are fewer than 1,000 WR stars in the rest of the Local Group galaxies, with around 166 known in the Magellanic Clouds, 206 in the Triangulum Galaxy, and 154 in the Andromeda Galaxy.
Outside the local group, whole galaxy surveys have found thousands more WR stars and candidates. For example, in the M101 Group, over a thousand potential WR stars have been detected, from magnitude 21 to 25, and astronomers hope to eventually catalog over ten thousand. These stars are expected to be particularly common in the Wolf–Rayet galaxies named after them and in starburst galaxies.
Their characteristic emission lines are formed in the extended and dense high-velocity wind region enveloping the very hot stellar photosphere, which produces a flood of UV radiation that causes fluorescence in the line-forming wind region. This ejection process uncovers in succession, first the nitrogen-rich products of CNO cycle burning of hydrogen (WN stars), and later the carbon-rich layer due to He burning (WC and WO-type stars).
It can be seen that the WNh stars are completely different objects from the WN stars without hydrogen. Despite the similar spectra, they are much more massive, much larger, and some of the most luminous stars known. They have been detected as early as WN5h in the Magellanic Clouds. The nitrogen seen in the spectrum of WNh stars is still the product of CNO cycle fusion in the core, but it appears at the surface of the most massive stars due to rotational and convectional mixing while still in the core hydrogen burning phase, rather than after the outer envelope is lost during core helium fusion.
Some Wolf–Rayet stars of the carbon sequence ("WC"), especially those belonging to the latest types, are noticeable due to their production of dust. Usually this takes place on those belonging to binary systems as a product of the collision of the stellar winds forming the pair, as is the case of the famous binary WR 104; however this process occurs on single ones too.
A few – roughly 10% – of the central stars of planetary nebulae, despite their much lower masses – typically ~0.6 M☉ – are also observationally of the WR-type; i.e. they show emission line spectra with broad lines from helium, carbon and oxygen. Denoted [WR], they are much older objects descended from evolved low-mass stars and are closely related to white dwarfs, rather than to the very young, very massive population I stars that comprise the bulk of the WR class. These are now generally excluded from the class denoted as Wolf–Rayet stars, or referred to as Wolf–Rayet-type stars.
Metallicity
The numbers and properties of Wolf–Rayet stars vary with the chemical composition of their progenitor stars. A primary driver of this difference is the rate of mass loss at different levels of metallicity. Higher metallicity leads to high mass loss, which affects the evolution of massive stars and also the properties of Wolf–Rayet stars. Higher levels of mass loss cause stars to lose their outer layers before an iron core develops and collapses, so that the more massive red supergiants evolve back to hotter temperatures before exploding as a supernova, and the most massive stars never become red supergiants. In the Wolf–Rayet stage, higher mass loss leads to stronger depletion of the layers outside the convective core, lower hydrogen surface abundances and more rapid stripping of helium to produce a WC spectrum.
These trends can be observed in the various galaxies of the local group, where metallicity varies from near-solar levels in the Milky Way, somewhat lower in M31, lower still in the Large Magellanic Cloud, and much lower in the Small Magellanic Cloud. Strong metallicity variations are seen across individual galaxies, with M33 and the Milky Way showing higher metallicities closer to the centre, and M31 showing higher metallicity in the disk than in the halo. Thus the SMC is seen to have few WR stars compared to its stellar formation rate and no WC stars at all (one star has a WO spectral type), the Milky Way has roughly equal numbers of WN and WC stars and a large total number of WR stars, and the other main galaxies have somewhat fewer WR stars and more WN than WC types. LMC, and especially SMC, Wolf–Rayets have weaker emission and a tendency to higher atmospheric hydrogen fractions. SMC WR stars almost universally show some hydrogen and even absorption lines even at the earliest spectral types, due to weaker winds not entirely masking the photosphere.
The maximum mass of a main-sequence star that can evolve through a red supergiant phase and back to a WNL star is calculated to be around in the Milky Way, in the LMC, and in the SMC. The more evolved WNE and WC stages are only reached by stars with an initial mass over at near-solar metallicity, in the LMC. Normal single star evolution is not expected to produce any WNE or WC stars at SMC metallicity.
Rotation
Mass loss is influenced by a star's rotation rate, especially strongly at low metallicity. Fast rotation contributes to mixing of core fusion products through the rest of the star, enhancing surface abundances of heavy elements, and driving mass loss. Rotation causes stars to remain on the main sequence longer than non-rotating stars, evolve more quickly away from the red supergiant phase, or even evolve directly from the main sequence to hotter temperatures for very high masses, high metallicity or very rapid rotation.
Stellar mass loss produces a loss of angular momentum and this quickly brakes the rotation of massive stars. Very massive stars at near-solar metallicity should be braked almost to a standstill while still on the main sequence, while at SMC metallicity they can continue to rotate rapidly even at the highest observed masses. Rapid rotation of massive stars may account for the unexpected properties and numbers of SMC WR stars, for example their relatively high temperatures and luminosities.
Binaries
Massive stars in binary systems can develop into Wolf–Rayet stars due to stripping by a companion rather than inherent mass loss due to a stellar wind. This process is relatively insensitive to the metallicity or rotation of the individual stars and is expected to produce a consistent set of WR stars across all the local group galaxies. As a result, the fraction of WR stars produced through the binary channel, and therefore the number of WR stars observed to be in binaries, should be higher in low metallicity environments. Calculations suggest that the binary fraction of WR stars observed in the SMC should be as high as 98%, although less than half are actually observed to have a massive companion. The binary fraction in the Milky Way is around 20%, in line with theoretical calculations.
Nebulae
A significant proportion of WR stars are surrounded by nebulosity associated directly with the star, not just the normal background nebulosity associated with any massive star forming region, and not a planetary nebula formed by a post-AGB star. The nebulosity presents a variety of forms and classification has been difficult. Many were originally catalogued as planetary nebulae and sometimes only a careful multi-wavelength study can distinguish a planetary nebula around a low mass post-AGB star from a similarly shaped nebula around a more massive core helium-burning star.
Wolf–Rayet galaxies
A Wolf–Rayet galaxy is a type of starburst galaxy where a sufficient number of WR stars exist that their characteristic emission line spectra become visible in the overall spectrum of the galaxy. Specifically a broad emission feature due to the 468.6 nm He and nearby spectral lines is the defining characteristic of a Wolf–Rayet galaxy. The relatively short lifetime of WR stars means that the starbursts in such galaxies must have occurred within the last few million years, and must have lasted less than a million years or else the WR emission would be swamped by large numbers of other luminous stars.
Evolution
Theories about how WR stars form, develop, and die have been slow to form compared to the explanation of less extreme stellar evolution. They are rare, distant, and often obscured, and even into the 21st century many aspects of their lives are unclear.
History
Although Wolf–Rayet stars have been clearly identified as an unusual and distinctive class of stars since the 19th century,
the nature of these stars was uncertain until towards the end of the 20th century. Before the 1960s, even the classification of WR stars was highly uncertain, and their nature and evolution was essentially unknown. The very similar appearance of the central stars of planetary nebulae (CSPNe) and the much more luminous classical WR stars contributed to the uncertainty.
By about 1960, the distinction between CSPNe and massive luminous classical WR stars was more clear. Studies showed that they were small dense stars surrounded by extensive circumstellar material, but not yet clear whether the material was expelled from the star or contracting onto it.
The unusual abundances of nitrogen, carbon, and oxygen, as well as the lack of hydrogen, were recognised, but the reasons remained obscure.
It was recognised that WR stars were very young and very rare, but it was still open to debate whether they were evolving towards or away from the main sequence.
By the 1980s, WR stars were accepted as the descendants of massive OB stars, although their exact evolutionary state in relation to the main sequence and other evolved massive stars was still unknown.
Theories that the preponderance of WR stars in massive binaries and their lack of hydrogen could be due to gravitational stripping had been largely ignored or abandoned.
WR stars were being proposed as possible progenitors of supernovae, and particularly the newly-discovered type Ib supernovae, lacking hydrogen but apparently associated with young massive stars.
By the start of the 21st century, WR stars were largely accepted as massive stars that had exhausted their core hydrogen, left the main sequence, and expelled most of their atmospheres, leaving behind a small hot core of helium and heavier fusion products.
Current models
Most WR stars, the classical population I type, are now understood as being a natural stage in the evolution of the most massive stars (not counting the less common planetary nebula central stars), either after a period as a red supergiant, after a period as a blue supergiant, or directly from the most massive main-sequence stars. Only the lower mass red supergiants are expected to explode as a supernova at that stage, while more massive red supergiants progress back to hotter temperatures as they expel their atmospheres. Some explode while at the yellow hypergiant or LBV stage, but many become Wolf–Rayet stars.
They have lost or burnt almost all of their hydrogen and are now fusing helium in their cores, or heavier elements for a very brief period at the end of their lives.
Massive main-sequence stars create a very hot core which fuses hydrogen very rapidly via the CNO process and results in strong convection throughout the whole star. This causes mixing of helium to the surface, a process that is enhanced by rotation, possibly by differential rotation where the core is spun up to a faster rotation than the surface. Such stars also show nitrogen enhancement at the surface at a very young age, caused by changes in the proportions of carbon and nitrogen due to the CNO cycle. The enhancement of heavy elements in the atmosphere, as well as increases in luminosity, create strong stellar winds which are the source of the emission line spectra. These stars develop an Of spectrum, Of* if they are sufficiently hot, which develops into a WNh spectrum as the stellar winds increase further. This explains the high mass and luminosity of the WNh stars, which are still burning hydrogen at the core and have lost little of their initial mass. These will eventually expand into blue supergiants (LBVs?) as hydrogen at the core becomes depleted, or if mixing is efficient enough (e.g. through rapid rotation) they may progress directly to WN stars without hydrogen.
WR stars are likely to end their lives violently rather than fade away to a white dwarf. Thus every star with an initial mass more than about 9 times the Sun would inevitably result in a supernova explosion (with the exception of direct collapse), many of them from the WR stage.
A simple progression of WR stars from low to hot temperatures, resulting finally in WO-type stars, is not supported by observation. WO-type stars are extremely rare and all the known examples are more luminous and more massive than the relatively common WC stars. Alternative theories suggest either that the WO-type stars are only formed from the most massive main-sequence stars, and/or that they form an extremely short-lived end stage of just a few thousand years before exploding, with the WC phase corresponding to the core helium burning phase and the WO phase to nuclear burning stages beyond. It is still unclear whether the WO spectrum is purely the result of ionisation effects at very high temperature, reflects an actual chemical abundance difference, or if both effects occur to varying degrees.
Key:
O: O-type main-sequence star
Of: evolved O-type showing N and He emission
BSG: blue supergiant
RSG: red supergiant
YHG: yellow hypergiant
LBV: luminous blue variable
WNh: WN plus hydrogen lines
WNL: "late" WN-class Wolf–Rayet star (about WN6 to WN11)
WNE: "early" WN-class Wolf–Rayet star (about WN2 to WN6)
WN/WC: Transitional (transitioning from WN to WC) Wolf–Rayet star (may be WN#/WCE or WC#/WN)
WC: WC-class Wolf–Rayet star
WO: WO-class Wolf–Rayet star
Wolf–Rayet stars form from massive stars, although the evolved population I stars have lost half or more of their initial masses by the time they show a WR appearance. For example, γ2 Velorum A currently has a mass around 9 times the Sun, but began with a mass at least 40 times the Sun. High-mass stars are very rare, both because they form less often and because they have short lives. This means that Wolf–Rayet stars themselves are extremely rare because they only form from the most massive main-sequence stars and because they are a relatively short-lived phase in the lives of those stars. This also explains why type Ib/c supernovae are less common than type II, since they result from higher-mass stars.
WNh stars, spectroscopically similar but actually a much less evolved star which has only just started to expel its atmosphere, are an exception and still retain much of their initial mass. The most massive stars currently known are all WNh stars rather than O-type main-sequence stars, an expected situation because such stars show helium and nitrogen at the surface only a few thousand years after they form, possibly before they become visible through the surrounding gas cloud. An alternative explanation is that these stars are so massive that they could not form as normal main-sequence stars, instead being the result of mergers of less extreme stars.
The difficulties of modelling the observed numbers and types of Wolf–Rayet stars through single star evolution have led to theories that they form through binary interactions which could accelerate loss of the outer layers of a star through mass exchange. WR 122 is a potential example that has a flat disk of gas encircling the star, almost 2 trillion miles wide, and may have a companion star that stripped its outer envelope.
Supernovae
It is widely suspected that many type Ib and type Ic supernova progenitors are WR stars, although no conclusive identification has been made of such a progenitor.
Type Ib supernovae lack hydrogen lines in their spectra. The more common type Ic supernovae lack both hydrogen and helium lines in their spectra. The expected progenitors for such supernova are massive stars that respectively lack hydrogen in their outer layers, or lack both hydrogen and helium. WR stars are just such objects. All WR stars lack hydrogen and in some WR stars, most notably the WO group, helium is also strongly depleted. WR stars are expected to experience core collapse when they have generated an iron core, and resulting supernova explosions would be of type Ib or Ic. In some cases it is possible that direct collapse of the core to a black hole would not produce a visible explosion.
WR stars are very luminous due to their high temperatures but not visually bright, especially the hottest examples that are expected to make up most supernova progenitors. Theory suggests that the progenitors of type Ibc supernovae observed to date would not be bright enough to be detected, although they place constraints on the properties of those progenitors. A possible progenitor star which has disappeared at the location of supernova iPTF13bvn may be a single WR star, although other analyses favour a less massive binary system with a stripped star or helium giant. The only other possible WR supernova progenitor is for SN 2017ein, and again it is uncertain whether the progenitor is a single massive WR star or binary system.
In 2022 astronomers from the Gran Telescopio Canarias reported the first supernova explosion of a Wolf–Rayet star. SN 2019hgp was a type Icn supernova and is also the first in which the element neon has been detected.
Examples
By far the most visible example of a Wolf–Rayet star is γ2 Velorum (WR 11), which is a bright naked eye star for those located south of 40 degrees northern latitude, although most of the light comes from an O7.5 giant companion. Due to the exotic nature of its spectrum (bright emission lines in lieu of dark absorption lines) it is dubbed the "Spectral Gem of the Southern Skies". The only other Wolf–Rayet star brighter than magnitude 6 is θ Muscae (WR 48), a triple star with two O class companions. Both are WC stars. The "ex" WR star WR 79a (HR 6272) is brighter than magnitude 6 but is now considered to be a peculiar O8 supergiant with strong emission. The next brightest at magnitude 6.4 is WR 22, a massive binary with a WN7h primary.
The most massive and most luminous star currently known, R136a1, is also a Wolf–Rayet star of the WNh type that is still fusing hydrogen in its core. This type of star, which includes many of the most luminous and most massive stars, is very young and usually found only in the centre of the densest star clusters. Occasionally a runaway WNh star such as VFTS 682 is found outside such clusters, probably having been ejected from a multiple system or by interaction with other stars.
An example of a triple star system containing a Wolf–Rayet binary is Apep. It releases huge amounts of carbon dust driven by their extreme stellar winds. As the two stars orbit one another, the dust gets wrapped into a glowing sooty tail.
All of the very hottest non-degenerate stars (the hottest few) are Wolf–Rayet stars, the hottest of which being WR 102, which seems to be as hot as 210,000 K, followed by WR 142 which is around 200,000 K in temperature. LMC195-1, located in the Large Magellanic Cloud, should have a similar temperature, but at the moment this temperature is unknown.
HD 45166 has been described as the most magnetic massive star known and as the first magnetic known Wolf-Rayet star.
Only a minority of planetary nebulae have WR type central stars, but a considerable number of well-known planetary nebulae do have them.
| Physical sciences | Stellar astronomy | null |
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