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59038 | https://en.wikipedia.org/wiki/Fourier%20series | Fourier series | A Fourier series () is an expansion of a periodic function into a sum of trigonometric functions. The Fourier series is an example of a trigonometric series. By expressing a function as a sum of sines and cosines, many problems involving the function become easier to analyze because trigonometric functions are well understood. For example, Fourier series were first used by Joseph Fourier to find solutions to the heat equation. This application is possible because the derivatives of trigonometric functions fall into simple patterns. Fourier series cannot be used to approximate arbitrary functions, because most functions have infinitely many terms in their Fourier series, and the series do not always converge. Well-behaved functions, for example smooth functions, have Fourier series that converge to the original function. The coefficients of the Fourier series are determined by integrals of the function multiplied by trigonometric functions, described in Fourier series§Definition.
The study of the convergence of Fourier series focus on the behaviors of the partial sums, which means studying the behavior of the sum as more and more terms from the series are summed. The figures below illustrate some partial Fourier series results for the components of a square wave.
Fourier series are closely related to the Fourier transform, a more general tool that can even find the frequency information for functions that are not periodic. Periodic functions can be identified with functions on a circle; for this reason Fourier series are the subject of Fourier analysis on the circle group, denoted by or . The Fourier transform is also part of Fourier analysis, but is defined for functions on .
Since Fourier's time, many different approaches to defining and understanding the concept of Fourier series have been discovered, all of which are consistent with one another, but each of which emphasizes different aspects of the topic. Some of the more powerful and elegant approaches are based on mathematical ideas and tools that were not available in Fourier's time. Fourier originally defined the Fourier series for real-valued functions of real arguments, and used the sine and cosine functions in the decomposition. Many other Fourier-related transforms have since been defined, extending his initial idea to many applications and birthing an area of mathematics called Fourier analysis.
History
The Fourier series is named in honor of Jean-Baptiste Joseph Fourier (1768–1830), who made important contributions to the study of trigonometric series, after preliminary investigations by Leonhard Euler, Jean le Rond d'Alembert, and Daniel Bernoulli. Fourier introduced the series for the purpose of solving the heat equation in a metal plate, publishing his initial results in his 1807 Mémoire sur la propagation de la chaleur dans les corps solides (Treatise on the propagation of heat in solid bodies), and publishing his Théorie analytique de la chaleur (Analytical theory of heat) in 1822. The Mémoire introduced Fourier analysis, specifically Fourier series. Through Fourier's research the fact was established that an arbitrary (at first, continuous and later generalized to any piecewise-smooth) function can be represented by a trigonometric series. The first announcement of this great discovery was made by Fourier in 1807, before the French Academy. Early ideas of decomposing a periodic function into the sum of simple oscillating functions date back to the 3rd century BC, when ancient astronomers proposed an empiric model of planetary motions, based on deferents and epicycles.
The heat equation is a partial differential equation. Prior to Fourier's work, no solution to the heat equation was known in the general case, although particular solutions were known if the heat source behaved in a simple way, in particular, if the heat source was a sine or cosine wave. These simple solutions are now sometimes called eigensolutions. Fourier's idea was to model a complicated heat source as a superposition (or linear combination) of simple sine and cosine waves, and to write the solution as a superposition of the corresponding eigensolutions. This superposition or linear combination is called the Fourier series.
From a modern point of view, Fourier's results are somewhat informal, due to the lack of a precise notion of function and integral in the early nineteenth century. Later, Peter Gustav Lejeune Dirichlet and Bernhard Riemann expressed Fourier's results with greater precision and formality.
Although the original motivation was to solve the heat equation, it later became obvious that the same techniques could be applied to a wide array of mathematical and physical problems, and especially those involving linear differential equations with constant coefficients, for which the eigensolutions are sinusoids. The Fourier series has many such applications in electrical engineering, vibration analysis, acoustics, optics, signal processing, image processing, quantum mechanics, econometrics, shell theory, etc.
Beginnings
Joseph Fourier wrote
This immediately gives any coefficient ak of the trigonometric series for φ(y) for any function which has such an expansion. It works because if φ has such an expansion, then (under suitable convergence assumptions) the integral
can be carried out term-by-term. But all terms involving for vanish when integrated from −1 to 1, leaving only the term, which is 1.
In these few lines, which are close to the modern formalism used in Fourier series, Fourier revolutionized both mathematics and physics. Although similar trigonometric series were previously used by Euler, d'Alembert, Daniel Bernoulli and Gauss, Fourier believed that such trigonometric series could represent any arbitrary function. In what sense that is actually true is a somewhat subtle issue and the attempts over many years to clarify this idea have led to important discoveries in the theories of convergence, function spaces, and harmonic analysis.
When Fourier submitted a later competition essay in 1811, the committee (which included Lagrange, Laplace, Malus and Legendre, among others) concluded: ...the manner in which the author arrives at these equations is not exempt of difficulties and...his analysis to integrate them still leaves something to be desired on the score of generality and even rigour.
Fourier's motivation
The Fourier series expansion of the sawtooth function (below) looks more complicated than the simple formula , so it is not immediately apparent why one would need the Fourier series. While there are many applications, Fourier's motivation was in solving the heat equation. For example, consider a metal plate in the shape of a square whose sides measure meters, with coordinates . If there is no heat source within the plate, and if three of the four sides are held at 0 degrees Celsius, while the fourth side, given by , is maintained at the temperature gradient degrees Celsius, for in , then one can show that the stationary heat distribution (or the heat distribution after a long time has elapsed) is given by
Here, sinh is the hyperbolic sine function. This solution of the heat equation is obtained by multiplying each term of the equation from Analysis § Example by . While our example function seems to have a needlessly complicated Fourier series, the heat distribution is nontrivial. The function cannot be written as a closed-form expression. This method of solving the heat problem was made possible by Fourier's work.
Other applications
Another application is to solve the Basel problem by using Parseval's theorem. The example generalizes and one may compute ζ(2n), for any positive integer n.
Definition
The Fourier series of a complex-valued -periodic function , integrable over the interval on the real line, is defined as a trigonometric series of the form
such that the Fourier coefficients are complex numbers defined by the integral
The series does not necessarily converge (in the pointwise sense) and, even if it does, it is not necessarily equal to . Only when certain conditions are satisfied (e.g. if is continuously differentiable) does the Fourier series converge to , i.e.,
For functions satisfying the Dirichlet sufficiency conditions, pointwise convergence holds. However, these are not necessary conditions and there are many theorems about different types of convergence of Fourier series (e.g. uniform convergence or mean convergence). The definition naturally extends to the Fourier series of a (periodic) distribution (also called Fourier-Schwartz series). Then the Fourier series converges to in the distribution sense.
The process of determining the Fourier coefficients of a given function or signal is called analysis, while forming the associated trigonometric series (or its various approximations) is called synthesis.
Synthesis
A Fourier series can be written in several equivalent forms, shown here as the partial sums of the Fourier series of :
The harmonics are indexed by an integer, which is also the number of cycles the corresponding sinusoids make in interval . Therefore, the sinusoids have:
a wavelength equal to in the same units as .
a frequency equal to in the reciprocal units of .
These series can represent functions that are just a sum of one or more frequencies in the harmonic spectrum. In the limit , a trigonometric series can also represent the intermediate frequencies and/or non-sinusoidal functions because of the infinite number of terms.
Analysis
The coefficients can be given/assumed, such as a music synthesizer or time samples of a waveform. In the latter case, the exponential form of Fourier series synthesizes a discrete-time Fourier transform where variable represents frequency instead of time. In general, the coefficients are determined by analysis of a given function whose domain of definition is an interval of length .
The scale factor follows from substituting into and utilizing the orthogonality of the trigonometric system. The equivalence of and follows from Euler's formula
resulting in:
with being the mean value of on the interval . Conversely:
Example
Consider a sawtooth function:
In this case, the Fourier coefficients are given by
It can be shown that the Fourier series converges to at every point where is differentiable, and therefore:
When , the Fourier series converges to 0, which is the half-sum of the left- and right-limit of at . This is a particular instance of the Dirichlet theorem for Fourier series.
This example leads to a solution of the Basel problem.
Amplitude-phase form
If the function is real-valued then the Fourier series can also be represented as
where is the amplitude and is the phase shift of the harmonic.
The equivalence of and follows from the trigonometric identity:
which implies
are the rectangular coordinates of a vector with polar coordinates and given by
where is the argument of .
An example of determining the parameter for one value of is shown in Figure 2. It is the value of at the maximum correlation between and a cosine template, The blue graph is the cross-correlation function, also known as a matched filter:
Fortunately, it isn't necessary to evaluate this entire function, because its derivative is zero at the maximum:
Hence
Common notations
The notation is inadequate for discussing the Fourier coefficients of several different functions. Therefore, it is customarily replaced by a modified form of the function ( in this case), such as or and functional notation often replaces subscripting:
In engineering, particularly when the variable represents time, the coefficient sequence is called a frequency domain representation. Square brackets are often used to emphasize that the domain of this function is a discrete set of frequencies.
Another commonly used frequency domain representation uses the Fourier series coefficients to modulate a Dirac comb:
where represents a continuous frequency domain. When variable has units of seconds, has units of hertz. The "teeth" of the comb are spaced at multiples (i.e. harmonics) of , which is called the fundamental frequency. can be recovered from this representation by an inverse Fourier transform:
The constructed function is therefore commonly referred to as a Fourier transform, even though the Fourier integral of a periodic function is not convergent at the harmonic frequencies.
Table of common Fourier series
Some common pairs of periodic functions and their Fourier series coefficients are shown in the table below.
designates a periodic function with period
designate the Fourier series coefficients (sine-cosine form) of the periodic function
Table of basic transformation rules
This table shows some mathematical operations in the time domain and the corresponding effect in the Fourier series coefficients. Notation:
Complex conjugation is denoted by an asterisk.
designate -periodic functions or functions defined only for
designate the Fourier series coefficients (exponential form) of and
Properties
Symmetry relations
When the real and imaginary parts of a complex function are decomposed into their even and odd parts, there are four components, denoted below by the subscripts RE, RO, IE, and IO. And there is a one-to-one mapping between the four components of a complex time function and the four components of its complex frequency transform:
From this, various relationships are apparent, for example:
The transform of a real-valued function is the conjugate symmetric function Conversely, a conjugate symmetric transform implies a real-valued time-domain.
The transform of an imaginary-valued function is the conjugate antisymmetric function and the converse is true.
The transform of a conjugate symmetric function is the real-valued function and the converse is true.
The transform of a conjugate antisymmetric function is the imaginary-valued function and the converse is true.
Riemann–Lebesgue lemma
If is integrable, , and
Parseval's theorem
If belongs to (periodic over an interval of length ) then:
Plancherel's theorem
If are coefficients and then there is a unique function such that for every .
Convolution theorems
Given -periodic functions, and with Fourier series coefficients and
The pointwise product: is also -periodic, and its Fourier series coefficients are given by the discrete convolution of the and sequences:
The periodic convolution: is also -periodic, with Fourier series coefficients:
A doubly infinite sequence in is the sequence of Fourier coefficients of a function in if and only if it is a convolution of two sequences in . See
Derivative property
If is a 2-periodic function on which is times differentiable, and its derivative is continuous, then belongs to the function space .
If , then the Fourier coefficients of the derivative of can be expressed in terms of the Fourier coefficients of , via the formula In particular, since for any fixed we have as , it follows that tends to zero, i.e., the Fourier coefficients converge to zero faster than the power of .
Compact groups
One of the interesting properties of the Fourier transform which we have mentioned, is that it carries convolutions to pointwise products. If that is the property which we seek to preserve, one can produce Fourier series on any compact group. Typical examples include those classical groups that are compact. This generalizes the Fourier transform to all spaces of the form L2(G), where G is a compact group, in such a way that the Fourier transform carries convolutions to pointwise products. The Fourier series exists and converges in similar ways to the case.
An alternative extension to compact groups is the Peter–Weyl theorem, which proves results about representations of compact groups analogous to those about finite groups.
Riemannian manifolds
If the domain is not a group, then there is no intrinsically defined convolution. However, if is a compact Riemannian manifold, it has a Laplace–Beltrami operator. The Laplace–Beltrami operator is the differential operator that corresponds to Laplace operator for the Riemannian manifold . Then, by analogy, one can consider heat equations on . Since Fourier arrived at his basis by attempting to solve the heat equation, the natural generalization is to use the eigensolutions of the Laplace–Beltrami operator as a basis. This generalizes Fourier series to spaces of the type , where is a Riemannian manifold. The Fourier series converges in ways similar to the case. A typical example is to take to be the sphere with the usual metric, in which case the Fourier basis consists of spherical harmonics.
Locally compact Abelian groups
The generalization to compact groups discussed above does not generalize to noncompact, nonabelian groups. However, there is a straightforward generalization to Locally Compact Abelian (LCA) groups.
This generalizes the Fourier transform to or , where is an LCA group. If is compact, one also obtains a Fourier series, which converges similarly to the case, but if is noncompact, one obtains instead a Fourier integral. This generalization yields the usual Fourier transform when the underlying locally compact Abelian group is .
Extensions
Fourier-Stieltjes series
Let be a function of bounded variation defined on the closed inverval . The Fourier series whose coefficients are given by
is called the Fourier-Stieltjes series. The space of functions of bounded variation is a subspace of . As any defines a Radon measure (i.e. a locally finite Borel measure on ), this definition can be extended as follows.
Consider the space of all finite Borel measures on the real line; as such . If there is a measure such that the Fourier-Stieltjes coefficients are given by
then the series is called a Fourier-Stieltjes series. Likewise, the function , where , is called a Fourier-Stieltjes transform.
The question whether or not exists for a given sequence of forms the basis of the trigonometric moment problem.
Furthermore, is a strict subspace of the space of (tempered) distributions , i.e., . If the Fourier coefficients are determined by a distribution then the series is described as a Fourier-Schwartz series. Contrary to the Fourier-Stieltjes series, deciding whether a given series is a Fourier series or a Fourier-Schwartz series is relatively trivial due to the characteristics of its dual space; the Schwartz space .
Fourier series on a square
We can also define the Fourier series for functions of two variables and in the square :
Aside from being useful for solving partial differential equations such as the heat equation, one notable application of Fourier series on the square is in image compression. In particular, the JPEG image compression standard uses the two-dimensional discrete cosine transform, a discrete form of the Fourier cosine transform, which uses only cosine as the basis function.
For two-dimensional arrays with a staggered appearance, half of the Fourier series coefficients disappear, due to additional symmetry.
Fourier series of Bravais-lattice-periodic-function
A three-dimensional Bravais lattice is defined as the set of vectors of the form:
where are integers and are three linearly independent vectors. Assuming we have some function, , such that it obeys the condition of periodicity for any Bravais lattice vector , , we could make a Fourier series of it. This kind of function can be, for example, the effective potential that one electron "feels" inside a periodic crystal. It is useful to make the Fourier series of the potential when applying Bloch's theorem. First, we may write any arbitrary position vector in the coordinate-system of the lattice:
where meaning that is defined to be the magnitude of , so is the unit vector directed along .
Thus we can define a new function,
This new function, , is now a function of three-variables, each of which has periodicity , , and respectively:
This enables us to build up a set of Fourier coefficients, each being indexed by three independent integers . In what follows, we use function notation to denote these coefficients, where previously we used subscripts. If we write a series for on the interval for , we can define the following:
And then we can write:
Further defining:
We can write once again as:
Finally applying the same for the third coordinate, we define:
We write as:
Re-arranging:
Now, every reciprocal lattice vector can be written (but does not mean that it is the only way of writing) as , where are integers and are reciprocal lattice vectors to satisfy ( for , and for ). Then for any arbitrary reciprocal lattice vector and arbitrary position vector in the original Bravais lattice space, their scalar product is:
So it is clear that in our expansion of , the sum is actually over reciprocal lattice vectors:
where
Assuming
we can solve this system of three linear equations for , , and in terms of , and in order to calculate the volume element in the original rectangular coordinate system. Once we have , , and in terms of , and , we can calculate the Jacobian determinant:
which after some calculation and applying some non-trivial cross-product identities can be shown to be equal to:
(it may be advantageous for the sake of simplifying calculations, to work in such a rectangular coordinate system, in which it just so happens that is parallel to the x axis, lies in the xy-plane, and has components of all three axes). The denominator is exactly the volume of the primitive unit cell which is enclosed by the three primitive-vectors , and . In particular, we now know that
We can write now as an integral with the traditional coordinate system over the volume of the primitive cell, instead of with the , and variables:
writing for the volume element ; and where is the primitive unit cell, thus, is the volume of the primitive unit cell.
Hilbert space
As the trigonometric series is a special class of orthogonal system, Fourier series can naturally be defined in the context of Hilbert spaces. For example, the space of square-integrable functions on forms the Hilbert space . Its inner product, defined for any two elements and , is given by:
This space is equipped with the orthonormal basis .
Then the (generalized) Fourier series expansion of , given by
can be written as
The sine-cosine form follows in a similar fashion. Indeed, the sines and cosines form an orthogonal set:
(where δmn is the Kronecker delta), and
Hence, the set
also forms an orthonormal basis for . The density of their span is a consequence of the Stone–Weierstrass theorem, but follows also from the properties of classical kernels like the Fejér kernel.
Fourier theorem proving convergence of Fourier series
In engineering, the Fourier series is generally assumed to converge except at jump discontinuities since the functions encountered in engineering are usually better-behaved than those in other disciplines. In particular, if is continuous and the derivative of (which may not exist everywhere) is square integrable, then the Fourier series of converges absolutely and uniformly to . If a function is square-integrable on the interval , then the Fourier series converges to the function almost everywhere. It is possible to define Fourier coefficients for more general functions or distributions, in which case pointwise convergence often fails, and convergence in norm or weak convergence is usually studied.
The theorems proving that a Fourier series is a valid representation of any periodic function (that satisfies the Dirichlet conditions), and informal variations of them that don't specify the convergence conditions, are sometimes referred to generically as Fourier's theorem or the Fourier theorem.
Least squares property
The earlier :
is a trigonometric polynomial of degree that can be generally expressed as:
Parseval's theorem implies that:
Convergence theorems
Because of the least squares property, and because of the completeness of the Fourier basis, we obtain an elementary convergence result.
If is continuously differentiable, then is the Fourier coefficient of the first derivative . Since is continuous, and therefore bounded, it is square-integrable and its Fourier coefficients are square-summable. Then, by the Cauchy–Schwarz inequality,
This means that is absolutely summable. The sum of this series is a continuous function, equal to , since the Fourier series converges in to :
This result can be proven easily if is further assumed to be , since in that case tends to zero as . More generally, the Fourier series is absolutely summable, thus converges uniformly to , provided that satisfies a Hölder condition of order . In the absolutely summable case, the inequality:
proves uniform convergence.
Many other results concerning the convergence of Fourier series are known, ranging from the moderately simple result that the series converges at if is differentiable at , to more sophisticated results such as Carleson's theorem which states that the Fourier series of an function converges almost everywhere.
Divergence
Since Fourier series have such good convergence properties, many are often surprised by some of the negative results. For example, the Fourier series of a continuous T-periodic function need not converge pointwise. The uniform boundedness principle yields a simple non-constructive proof of this fact.
In 1922, Andrey Kolmogorov published an article titled Une série de Fourier-Lebesgue divergente presque partout in which he gave an example of a Lebesgue-integrable function whose Fourier series diverges almost everywhere. He later constructed an example of an integrable function whose Fourier series diverges everywhere.
It is possible to give explicit examples of a continuous function whose Fourier series diverges at 0: for instance, the even and 2π-periodic function f defined for all x in [0,π] by
Because the function is even the Fourier series contains only cosines:
The coefficients are:
As increases, the coefficients will be positive and increasing until they reach a value of about at for some and then become negative (starting with a value around ) and getting smaller, before starting a new such wave. At the Fourier series is simply the running sum of and this builds up to around
in the th wave before returning to around zero, showing that the series does not converge at zero but reaches higher and higher peaks. Note that though the function is continuous, it is not differentiable.
| Mathematics | Calculus and analysis | null |
59052 | https://en.wikipedia.org/wiki/Ensemble%20%28mathematical%20physics%29 | Ensemble (mathematical physics) | In physics, specifically statistical mechanics, an ensemble (also statistical ensemble) is an idealization consisting of a large number of virtual copies (sometimes infinitely many) of a system, considered all at once, each of which represents a possible state that the real system might be in. In other words, a statistical ensemble is a set of systems of particles used in statistical mechanics to describe a single
system. The concept of an ensemble was introduced by J. Willard Gibbs in 1902.
A thermodynamic ensemble is a specific variety of statistical ensemble that, among other properties, is in statistical equilibrium (defined below), and is used to derive the properties of thermodynamic systems from the laws of classical or quantum mechanics.
Physical considerations
The ensemble formalises the notion that an experimenter repeating an experiment again and again under the same macroscopic conditions, but unable to control the microscopic details, may expect to observe a range of different outcomes.
The notional size of ensembles in thermodynamics, statistical mechanics and quantum statistical mechanics can be very large, including every possible microscopic state the system could be in, consistent with its observed macroscopic properties. For many important physical cases, it is possible to calculate averages directly over the whole of the thermodynamic ensemble, to obtain explicit formulas for many of the thermodynamic quantities of interest, often in terms of the appropriate partition function.
The concept of an equilibrium or stationary ensemble is crucial to many applications of statistical ensembles. Although a mechanical system certainly evolves over time, the ensemble does not necessarily have to evolve. In fact, the ensemble will not evolve if it contains all past and future phases of the system. Such a statistical ensemble, one that does not change over time, is called stationary and can be said to be in statistical equilibrium.
Terminology
The word "ensemble" is also used for a smaller set of possibilities sampled from the full set of possible states. For example, a collection of walkers in a Markov chain Monte Carlo iteration is called an ensemble in some of the literature.
The term "ensemble" is often used in physics and the physics-influenced literature. In probability theory, the term probability space is more prevalent.
Main types
The study of thermodynamics is concerned with systems that appear to human perception to be "static" (despite the motion of their internal parts), and which can be described simply by a set of macroscopically observable variables. These systems can be described by statistical ensembles that depend on a few observable parameters, and which are in statistical equilibrium. Gibbs noted that different macroscopic constraints lead to different types of ensembles, with particular statistical characteristics.
"We may imagine a great number of systems of the same nature, but differing in the configurations and velocities which they have at a given instant, and differing in not merely infinitesimally, but it may be so as to embrace every conceivable combination of configuration and velocities..." J. W. Gibbs (1903)
Three important thermodynamic ensembles were defined by Gibbs:
Microcanonical ensemble (or NVE ensemble) —a statistical ensemble where the total energy of the system and the number of particles in the system are each fixed to particular values; each of the members of the ensemble are required to have the same total energy and particle number. The system must remain totally isolated (unable to exchange energy or particles with its environment) in order to stay in statistical equilibrium.
Canonical ensemble (or NVT ensemble)—a statistical ensemble where the energy is not known exactly but the number of particles is fixed. In place of the energy, the temperature is specified. The canonical ensemble is appropriate for describing a closed system which is in, or has been in, weak thermal contact with a heat bath. In order to be in statistical equilibrium, the system must remain totally closed (unable to exchange particles with its environment) and may come into weak thermal contact with other systems that are described by ensembles with the same temperature.
Grand canonical ensemble (or μVT ensemble)—a statistical ensemble where neither the energy nor particle number are fixed. In their place, the temperature and chemical potential are specified. The grand canonical ensemble is appropriate for describing an open system: one which is in, or has been in, weak contact with a reservoir (thermal contact, chemical contact, radiative contact, electrical contact, etc.). The ensemble remains in statistical equilibrium if the system comes into weak contact with other systems that are described by ensembles with the same temperature and chemical potential.
The calculations that can be made using each of these ensembles are explored further in their respective articles.
Other thermodynamic ensembles can be also defined, corresponding to different physical requirements, for which analogous formulae can often similarly be derived.
For example, in the reaction ensemble, particle number fluctuations are only allowed to occur according to the stoichiometry of the chemical reactions which are present in the system.
Equivalence
In thermodynamic limit all ensembles should produce identical observables due to Legendre transforms, deviations to this rule occurs under conditions that state-variables are non-convex, such as small molecular measurements.
Representations
The precise mathematical expression for a statistical ensemble has a distinct form depending on the type of mechanics under consideration (quantum or classical). In the classical case, the ensemble is a probability distribution over the microstates. In quantum mechanics, this notion, due to von Neumann, is a way of assigning a probability distribution over the results of each complete set of commuting observables.
In classical mechanics, the ensemble is instead written as a probability distribution in phase space; the microstates are the result of partitioning phase space into equal-sized units, although the size of these units can be chosen somewhat arbitrarily.
Requirements for representations
Putting aside for the moment the question of how statistical ensembles are generated operationally, we should be able to perform the following two operations on ensembles A, B of the same system:
Test whether A, B are statistically equivalent.
If p is a real number such that , then produce a new ensemble by probabilistic sampling from A with probability p and from B with probability .
Under certain conditions, therefore, equivalence classes of statistical ensembles have the structure of a convex set.
Quantum mechanical
A statistical ensemble in quantum mechanics (also known as a mixed state) is most often represented by a density matrix, denoted by . The density matrix provides a fully general tool that can incorporate both quantum uncertainties (present even if the state of the system were completely known) and classical uncertainties (due to a lack of knowledge) in a unified manner. Any physical observable in quantum mechanics can be written as an operator, . The expectation value of this operator on the statistical ensemble is given by the following trace:
This can be used to evaluate averages (operator ), variances (using operator ), covariances (using operator ), etc. The density matrix must always have a trace of 1: (this essentially is the condition that the probabilities must add up to one).
In general, the ensemble evolves over time according to the von Neumann equation.
Equilibrium ensembles (those that do not evolve over time, ) can be written solely as a function of conserved variables. For example, the microcanonical ensemble and canonical ensemble are strictly functions of the total energy, which is measured by the total energy operator (Hamiltonian). The grand canonical ensemble is additionally a function of the particle number, measured by the total particle number operator . Such equilibrium ensembles are a diagonal matrix in the orthogonal basis of states that simultaneously diagonalize each conserved variable. In bra–ket notation, the density matrix is
where the , indexed by , are the elements of a complete and orthogonal basis. (Note that in other bases, the density matrix is not necessarily diagonal.)
Classical mechanical
In classical mechanics, an ensemble is represented by a probability density function defined over the system's phase space. While an individual system evolves according to Hamilton's equations, the density function (the ensemble) evolves over time according to Liouville's equation.
In a mechanical system with a defined number of parts, the phase space has generalized coordinates called , and associated canonical momenta called . The ensemble is then represented by a joint probability density function .
If the number of parts in the system is allowed to vary among the systems in the ensemble (as in a grand ensemble where the number of particles is a random quantity), then it is a probability distribution over an extended phase space that includes further variables such as particle numbers (first kind of particle), (second kind of particle), and so on up to (the last kind of particle; is how many different kinds of particles there are). The ensemble is then represented by a joint probability density function . The number of coordinates varies with the numbers of particles.
Any mechanical quantity can be written as a function of the system's phase. The expectation value of any such quantity is given by an integral over the entire phase space of this quantity weighted by :
The condition of probability normalization applies, requiring
Phase space is a continuous space containing an infinite number of distinct physical states within any small region. In order to connect the probability density in phase space to a probability distribution over microstates, it is necessary to somehow partition the phase space into blocks that are distributed representing the different states of the system in a fair way. It turns out that the correct way to do this simply results in equal-sized blocks of canonical phase space, and so a microstate in classical mechanics is an extended region in the phase space of canonical coordinates that has a particular volume. In particular, the probability density function in phase space, , is related to the probability distribution over microstates, by a factor
where
is an arbitrary but predetermined constant with the units of , setting the extent of the microstate and providing correct dimensions to .
is an overcounting correction factor (see below), generally dependent on the number of particles and similar concerns.
Since can be chosen arbitrarily, the notional size of a microstate is also arbitrary. Still, the value of influences the offsets of quantities such as entropy and chemical potential, and so it is important to be consistent with the value of when comparing different systems.
Correcting overcounting in phase space
Typically, the phase space contains duplicates of the same physical state in multiple distinct locations. This is a consequence of the way that a physical state is encoded into mathematical coordinates; the simplest choice of coordinate system often allows a state to be encoded in multiple ways. An example of this is a gas of identical particles whose state is written in terms of the particles' individual positions and momenta: when two particles are exchanged, the resulting point in phase space is different, and yet it corresponds to an identical physical state of the system. It is important in statistical mechanics (a theory about physical states) to recognize that the phase space is just a mathematical construction, and to not naively overcount actual physical states when integrating over phase space. Overcounting can cause serious problems:
Dependence of derived quantities (such as entropy and chemical potential) on the choice of coordinate system, since one coordinate system might show more or less overcounting than another.
Erroneous conclusions that are inconsistent with physical experience, as in the mixing paradox.
Foundational issues in defining the chemical potential and the grand canonical ensemble.
It is in general difficult to find a coordinate system that uniquely encodes each physical state. As a result, it is usually necessary to use a coordinate system with multiple copies of each state, and then to recognize and remove the overcounting.
A crude way to remove the overcounting would be to manually define a subregion of phase space that includes each physical state only once and then exclude all other parts of phase space. In a gas, for example, one could include only those phases where the particles' coordinates are sorted in ascending order. While this would solve the problem, the resulting integral over phase space would be tedious to perform due to its unusual boundary shape. (In this case, the factor introduced above would be set to , and the integral would be restricted to the selected subregion of phase space.)
A simpler way to correct the overcounting is to integrate over all of phase space but to reduce the weight of each phase in order to exactly compensate the overcounting. This is accomplished by the factor introduced above, which is a whole number that represents how many ways a physical state can be represented in phase space. Its value does not vary with the continuous canonical coordinates, so overcounting can be corrected simply by integrating over the full range of canonical coordinates, then dividing the result by the overcounting factor. However, does vary strongly with discrete variables such as numbers of particles, and so it must be applied before summing over particle numbers.
As mentioned above, the classic example of this overcounting is for a fluid system containing various kinds of particles, where any two particles of the same kind are indistinguishable and exchangeable. When the state is written in terms of the particles' individual positions and momenta, then the overcounting related to the exchange of identical particles is corrected by using
This is known as "correct Boltzmann counting".
Ensembles in statistics
The formulation of statistical ensembles used in physics has now been widely adopted in other fields, in part because it has been recognized that the canonical ensemble or Gibbs measure serves to maximize the entropy of a system, subject to a set of constraints: this is the principle of maximum entropy. This principle has now been widely applied to problems in linguistics, robotics, and the like.
In addition, statistical ensembles in physics are often built on a principle of locality: that all interactions are only between neighboring atoms or nearby molecules. Thus, for example, lattice models, such as the Ising model, model ferromagnetic materials by means of nearest-neighbor interactions between spins. The statistical formulation of the principle of locality is now seen to be a form of the Markov property in the broad sense; nearest neighbors are now Markov blankets. Thus, the general notion of a statistical ensemble with nearest-neighbor interactions leads to Markov random fields, which again find broad applicability; for example in Hopfield networks.
Ensemble average
In statistical mechanics, the ensemble average is defined as the mean of a quantity that is a function of the microstate of a system, according to the distribution of the system on its micro-states in this ensemble.
Since the ensemble average is dependent on the ensemble chosen, its mathematical expression varies from ensemble to ensemble. However, the mean obtained for a given physical quantity does not depend on the ensemble chosen at the thermodynamic limit.
The grand canonical ensemble is an example of an open system.
Classical statistical mechanics
For a classical system in thermal equilibrium with its environment, the ensemble average takes the form of an integral over the phase space of the system:
where
is the ensemble average of the system property A,
is , known as thermodynamic beta,
H is the Hamiltonian of the classical system in terms of the set of coordinates and their conjugate generalized momenta ,
is the volume element of the classical phase space of interest.
The denominator in this expression is known as the partition function and is denoted by the letter Z.
Quantum statistical mechanics
In quantum statistical mechanics, for a quantum system in thermal equilibrium with its environment, the weighted average takes the form of a sum over quantum energy states, rather than a continuous integral:
Canonical ensemble average
The generalized version of the partition function provides the complete framework for working with ensemble averages in thermodynamics, information theory, statistical mechanics and quantum mechanics.
The microcanonical ensemble represents an isolated system in which energy (E), volume (V) and the number of particles (N) are all constant. The canonical ensemble represents a closed system which can exchange energy (E) with its surroundings (usually a heat bath), but the volume (V) and the number of particles (N) are all constant. The grand canonical ensemble represents an open system which can exchange energy (E) and particles (N) with its surroundings, but the volume (V) is kept constant.
Operational interpretation
In the discussion given so far, while rigorous, we have taken for granted that the notion of an ensemble is valid a priori, as is commonly done in physical context. What has not been shown is that the ensemble itself (not the consequent results) is a precisely defined object mathematically. For instance,
It is not clear where this very large set of systems exists (for example, is it a gas of particles inside a container?)
It is not clear how to physically generate an ensemble.
In this section, we attempt to partially answer this question.
Suppose we have a preparation procedure for a system in a physics lab: For example, the procedure might involve a physical apparatus and some protocols for manipulating the apparatus. As a result of this preparation procedure, some system is produced and maintained in isolation for some small period of time. By repeating this laboratory preparation procedure we obtain a sequence of systems X1, X2, ...,Xk, which in our mathematical idealization, we assume is an infinite sequence of systems. The systems are similar in that they were all produced in the same way. This infinite sequence is an ensemble.
In a laboratory setting, each one of these prepped systems might be used as input for one subsequent testing procedure. Again, the testing procedure involves a physical apparatus and some protocols; as a result of the testing procedure we obtain a yes or no answer. Given a testing procedure E applied to each prepared system, we obtain a sequence of values Meas (E, X1), Meas (E, X2), ..., Meas (E, Xk). Each one of these values is a 0 (or no) or a 1 (yes).
Assume the following time average exists:
For quantum mechanical systems, an important assumption made in the quantum logic approach to quantum mechanics is the identification of yes–no questions to the lattice of closed subspaces of a Hilbert space. With some additional technical assumptions one can then infer that states are given by density operators S so that:
We see this reflects the definition of quantum states in general: A quantum state is a mapping from the observables to their expectation values.
| Physical sciences | Statistical mechanics | Physics |
59166 | https://en.wikipedia.org/wiki/Cassowary | Cassowary | Cassowaries (; Biak: man suar ; ; Papuan: kasu weri ) are flightless birds of the genus Casuarius, in the order Casuariiformes. They are classified as ratites, flightless birds without a keel on their sternum bones. Cassowaries are native to the tropical forests of New Guinea (Western New Guinea and Papua New Guinea), the Moluccas (Seram and Aru Islands), and northeastern Australia.
Three cassowary species are extant. The most common, the southern cassowary, is the third-tallest and second-heaviest living bird, smaller only than the ostrich and emu. The other two species are the northern cassowary and the dwarf cassowary; the northern cassowary is the most recently discovered and the most threatened. A fourth, extinct, species is the pygmy cassowary.
Cassowaries are very wary of humans, but if provoked, they are capable of inflicting serious, even fatal, injuries. They are known to attack both dogs and people. The cassowary has often been labelled "the world's most dangerous bird", although in terms of recorded statistics, it pales in comparison to the common ostrich, which kills two to three humans per year in South Africa.
Taxonomy, systematics, and evolution
The genus Casuarius was erected by French scientist Mathurin Jacques Brisson in his published in 1760. The type species is the southern cassowary (Casuarius casuarius). The Swedish naturalist Carl Linnaeus had introduced the genus Casuarius in the sixth edition of his published in 1748, but Linnaeus dropped the genus in the important tenth edition of 1758 and put the southern cassowary together with the common ostrich and the greater rhea in the genus Struthio. As the publication date of Linnaeus's sixth edition was before the 1758 starting point of the International Commission on Zoological Nomenclature, Brisson, and not Linnaeus, is considered the authority for the genus.
Cassowaries (from cognate of several related languages spoken around the Moluccas and New Guinea) are part of the ratite group, which also includes the emu, rheas, ostriches, and kiwi, as well as the extinct moas and elephant birds. These species are recognised:
Most authorities consider the taxonomic classification above to be monotypic, but several subspecies of each have been described, and some of them have even been suggested as separate species, e.g., C. (b) papuanus. The taxonomic name C. (b) papuanus also may be in need of revision to Casuarius (bennetti) westermanni. Validation of these subspecies has proven difficult due to individual variations, age-related variations, the scarcity of specimens, the stability of specimens (the bright skin of the head and neck—the basis of describing several subspecies—fades in specimens), and the practice of trading live cassowaries for thousands of years, some of which are likely to have escaped or been deliberately introduced to regions away from their origin.
The evolutionary history of cassowaries, as of all ratites, is not well known. Genetic evidence suggests that their closest living relatives are emus, and that the dwarf cassowary is more closely related to the Northern Cassowary than either is to the Southern cassowary. A fossil species was reported from Australia, but for reasons of biogeography, this assignment is not certain, and it might belong to the prehistoric Emuarius, which was a genus of cassowary-like primitive emus.
Description
Typically, all cassowaries are shy birds that are found in the deep forest. They are adept at disappearing long before a human knows they are there. The southern cassowary of the far north Queensland rain forests is not well studied, and the northern and dwarf cassowaries even less so. Females are larger and more brightly coloured than the males. Adult southern cassowaries are tall, although some females may reach , and weigh . However, it is not uncommon to see exceptionally large females topping the scales beyond , with the largest maximum recorded being a southern cassowary at and tall.
Hence, by technicality, all three species of cassowaries are considered as Asia's largest bird since the extinction of the Arabian ostrich. Moreover, not only is the cassowary Asia's largest bird, within New Guinea, the cassowary is the island's second largest terrestrial animal after the introduction of Cervidaes such as the rusa deer, chital, and fallow deer.
All cassowaries' feathers consist of a shaft and loose barbules. They do not have rectrices (tail feathers) or a preen gland. Cassowaries have small wings with five or six large remiges. These are reduced to stiff, keratinous quills, resembling porcupine quills, with no barbs. The furcula and coracoid are degenerate, and their palatal bones and sphenoid bones touch each other. These, along with their wedge-shaped body, are thought to be adaptations to ward off vines, thorns, and saw-edged leaves, allowing them to run quickly through the rainforest.
Unlike the majority of birds, cassowaries lack a tongue. Their beaks are pointed, sharp and robust but not serrated, which allows them to pick up fruit more easily than the short bills of an emu or an ostrich.
Cassowaries have three-toed feet with sharp claws. The inner (first) toe has a dagger-like claw that may be long. This claw is particularly fearsome, since cassowaries sometimes kick humans and other animals with their powerful legs. Cassowaries can run at up to through the dense forest and can jump up to . They are good swimmers, crossing wide rivers and swimming in the sea.
All three species have a keratinous, skin-covered casque on their heads that grows with age. The casque's shape and size, up to , is species-dependent. C. casuarius has the largest and C. bennetti the smallest (tricorn shape), with C. unappendiculatus having variations in between. Contrary to earlier findings, the hollow inside of the casque is spanned with fine fibres.
Several functions for the casque have been proposed. One is that they are a secondary sexual characteristic. Other suggested functions include batting through the underbrush, as a weapon in dominance disputes, or pushing aside leaf litter during foraging. The latter three are disputed by biologist Andrew Mack, whose personal observation suggests that the casque amplifies deep sounds. This is related to a discovery that at least the dwarf cassowary and southern cassowary produce very low-frequency sounds, which may aid in communication in dense rainforests. The "boom" vocalization that cassowaries produce is the lowest-frequency bird call known and is at the lower limit of human hearing. Recent study suggests that casque acts as a thermal radiator, offloading heat at high temperatures and restricting heat loss at low temperatures.
The average lifespan of wild cassowaries is approximately 18–20 years, with those held in captivity living up to 40 years.
Behaviour and ecology
Cassowaries are solitary birds except during courtship, egg-laying, and sometimes around ample food supplies. Males and females each maintain separate territories that overlap, of a size of approximately 3 square kilometres in one study. While females move among satellite territories of different males, they appear to remain within the same territories for most of their lives, mating with the same, or closely related, males over the course of their lives.
Courtship and pair-bonding rituals begin with the vibratory sounds broadcast by females. Males approach and run with their necks parallel to the ground while making dramatic movements of their heads, which accentuate the frontal neck region. The female approaches drumming slowly. The male crouches on the ground, and the female either steps on the male's back for a moment before crouching beside him in preparation for copulation, or she may attack. This is often the case with the females pursuing the males in ritualistic chasing behaviours that generally terminate in water. The male cassowary dives into water and submerges himself up to his upper neck and head. The female pursues him into the water, where he eventually drives her to the shallows, where she crouches making ritualistic motions of her head. The two may remain in copulation for extended periods of time. In some cases, another male may approach and run off the first male. He will climb onto her to copulate, as well.
Both male and female cassowaries do not tolerate the presence of others of the same sex, but females are more prone to fight than males, which will generally flee when encountering another male. While males and females may also be territorial and confrontational, this decreases during the mating season
Reproduction
The cassowary breeding season starts in May to June. Females lay three to eight large, bright green or pale green-blue eggs in each clutch into a heap of leaf litter prepared by the male. The eggs measure about – only ostrich and emu eggs are larger.
The male incubates those eggs for 50–52 days, removing or adding litter to regulate the temperature, then protects the chicks, which stay in the nest for about 9 months. He defends them fiercely against all potential predators, including humans. The young males later go off to find a territory of their own.
The female does not care for the eggs or the chicks, but rather moves on within her territory to lay eggs in the nests of several other males. Young cassowaries are brown and have buffy stripes. They are often kept as pets in native villages (in New Guinea), where they are permitted to roam like barnyard fowl until nearing maturity. Caged birds are regularly bereft of their fresh plumes.
Diet
Fruit from at least 26 plant families has been documented in the diet of cassowaries. Fruits from the laurel, podocarp, palm, wild grape, nightshade, and myrtle families are important items in the diet. The poisonous cassowary plum takes its name from the bird. The bird avoids the poisons of these fruits due to the presence of their incredibly short gastrointestinal tract, the shortest of all ratites in relation to their size. The cassowary's incredibly short and simple digestive tract leads to a short gut retention time which allow seeds to remain unharmed during the comparatively soft digestion process and allows them to consume fruits that contain toxins such as cyanogens. This short gut length also allows the birds to eat a wider variety of protein source, which is unsurprising given their omnivorous diet.
Where trees are dropping fruit, cassowaries come in and feed, with each bird defending a tree from others for a few days. They move on when the fruit is depleted. Fruit, even items as large as bananas and apples, is swallowed whole. Cassowaries are a keystone species of rain forests because they eat fallen fruit whole and distribute seeds across the jungle floor via excrement.
Adult and young cassowaries also practice coprophagia. As adult waste often contain half-digested fruit which still has nutritional value, so the birds would devour each other's as well as their own droppings.
In more urbanised areas, especially in Queensland, Australia, 'urbanised' cassowaries have adopted to also feed on picnic blankets, tables and baskets or backyard bird feeders and compost heaps, thereby consuming a wide range of non-natural and non-native foods as well. In fact, cassowaries are known to eat non-edible items—in one case, collection of urban cassowary droppings resulted in many unusual items. Outside of the skeletal remains of a honeyeater, researchers also found remains of a child’s coloured building blocks, various sized marbles and a very small plastic car that came from a cereal packet. In terms of roadkill, discarded fish was reported; another type of roadkill reported eaten by cassowaries is the bandicoot.
In captivity, cassowaries get the majority of their protein source from dog or monkey food. In fact, captive cassowaries consume almost of a protein source (such as dog food) in conjunction with of fruit a day, which results in 5% of their overall diet.
Role in seed dispersal and germination
Cassowaries feed on the fruit of several hundred rainforest species and usually pass viable seeds in large, dense scats. They are known to disperse seeds over distances greater than a kilometre, thus playing an important role in the ecosystem. Germination rates for seeds of the rare Australian rainforest tree Ryparosa were found to be much higher after passing through a cassowary's gut (92% versus 4%).
Threats
In its main home of New Guinea, cassowaries are the island's largest and most dominant and formidable bird, as well as being one of the largest terrestrial endemic animals in New Guinea. As such, adult cassowaries have no natural enemies other than humans (and even then, the birds are rarely hunted due to their reputation, speed, wariness and self-defence, with juveniles being preferred over adults for ceremonial purposes - on average, it is considered very fortunate for a human hunter to kill one in every five years). With regards to their relationship with the New Guinea singing dogs - one of Papua's only obligate terrestrial apex predator, with the other being the crocodile monitor - adult birds generally ignore them, with some even believing that the dogs take full advantage of the birds' foraging behaviour, as both species share and use the same feeding trail through the forests. It was believed that these dogs follow adult birds to catch small prey attracted to the dropped fruits on the rainforest floor. Nevertheless, there was a report from a native hunter of an exceptionally rare case of a singing dog attacking the dwarf species. The incident ended with the singing dog being disemboweled and ripped open by the bird. But generally speaking, both apex animals mutually keep their distance and avoid one another.
The same cannot be said with their chicks, however, as they are vulnerable to large pythons, monitor lizards, New Guinea singing dogs, and Papuan eagles. When threatened, it is known that cassowary chicks emit different vocalisation calls to indicate the specific threat, such as a hawk for example, before running underneath their father. Adult males aggressively defend their chicks. While adult males usually scare off or kill most predators, a chick will occasionally be separated in the chaos and become a potential target.
However, in the relic populations of north-eastern Australia, the cassowary population faces threats from vehicles, and are in danger of being outcompeted by wild boars, with their eggs being most vulnerable to boar predation. Their chicks also face dangers and predation from domesticated dogs, which results in a widespread decline in the Australian mainland. Because of such frequent inter-species conflicts, hunting dogs are one of the biggest enemies for cassowaries, and it is not unheard of for hunting dogs to accidentally kill cassowary chicks instead of feral pigs, with the dogs in-turn, being killed by the nearby adult rooster. Outside of threats from invasive species, the birds are also vulnerable to being unintentionally poisoned as well. It is unknown why the cassowary population in Australia is in decline, as the New Guinea population has dealt with introduced wild boars, dogs and feral cats for thousands of years longer with little to no impact on its population, suggesting that either the cassowaries of New Guinea had long adapted to human-introduced species or that the rich biodiversity of New Guinea allowed for additional niche partitioning.
As for eating the cassowary, it is supposed to be quite tough. Australian administrative officers stationed in New Guinea were advised that it "should be cooked with a stone in the pot: when the stone is ready to eat, so is the cassowary".
Distribution and habitat
Cassowaries are native to the humid rainforests of New Guinea, nearby smaller islands, East Nusa Tenggara, the Maluku Islands, and northeastern Australia. They do, however, venture out into palm scrub, grassland, savanna, and swamp forest. The wild population of cassowaries is threatened by deforestation, hunting, and habitat destruction. Human presence and agricultural activities have also contributed to the decline of their population in some areas. To protect this species, various conservation efforts have been carried out, including preserving natural habitat and enforcing regulations against illegal hunting. In Indonesia, cassowaries are predominantly found in the rainforests of Papua, particularly in lowland and montane areas. In addition, cassowaries are known to inhabit protected areas such as Wasur National Park in Merauke and Lorentz National Park, which is the largest national park in Southeast Asia, encompassing a vast range of ecosystems from coastal to alpine environments. These birds play a critical ecological role in seed dispersal, contributing to the regeneration of forests in these protected areas. They can also be easily spotted in some national parks such as Mellwraith Range National Park, Paluma Range National Park, and Jardine National Park in Australia.
Status and conservation
The southern cassowary is endangered in Queensland. Kofron and Chapman, when they assessed the decline of this species, found that of the former cassowary habitat, only 20–25% remains. Habitat loss and fragmentation is the primary cause of decline. They then studied 140 cases of cassowary mortality, and found that motor-vehicle strikes accounted for 55% of the deaths, and dog attacks produced another 18%. Remaining causes of death included hunting (five cases), entanglement in wire (one case), the removal of cassowaries that attacked humans (four cases), and natural causes (18 cases), including tuberculosis (four cases). The cause for 14 cases was indicated as "for unknown reasons".
Hand feeding cassowaries poses a significant threat to their survival because it lures them into suburban areas. There, the birds are more susceptible to encounters with vehicles and dogs. Contact with humans encourages cassowaries to take food from picnic tables. Feral pigs also are a significant threat to their survival. They destroy nests and eggs of cassowaries, but their worst effect is as competitors for food, which may be catastrophic for the cassowaries during lean times.
In February 2011, Cyclone Yasi destroyed a large area of cassowary habitat, endangering 200 of the birds – about 10% of the total Australian population.
The Mission Beach community in far north Queensland holds an annual Cassowary Festival in September, where funds are raised to map the bird's habitat.
Ironically, despite being a threatened species in Queensland, concerns also mount on cassowaries being itself, a potential invasive species on the island of Tasmania. According to the Department of Primary Industries, Parks, Water and Environment from Hobart, risk assessments on the cassowary as a potential invasive pest states that whilst the birds may have trouble establishing a stable population on the island, they would nonetheless, be considered a destructive element to Tasmania's ecological diversity and recommends strict imports on these birds. Reasons are many. The most notable are the birds' size. Cassowaries would automatically become the island's largest and most dominant terrestrial animal that could bully smaller animals in the same ecological niche. Frugivores such as the common brushtail possum, common ringtail possum, eastern pygmy possum and the little pygmy possum would be denied access to fruit of which they depend upon.
However, since Tasmania lack the same levels of fruit diversity as Queensland and New Guinea, assessments believe that the birds would adapt by also eating invertebrates and small vertebrates. This may lead to some competition with the island's endemic insectivores such as the eastern quoll, southern brown bandicoot, and eastern barred bandicoot.
In captivity
The cassowary has solitary habits and breeds less frequently in zoos than other ratites such as ostrich and emu. Unlike other ratites, it lives exclusively in tropical rainforest, and reproducing this habitat carefully is essential. Unlike the emu, which will live with other sympatric species, such as kangaroos, in "mixed Australian fauna" displays, the cassowary does not cohabit well among its own kind. Individual specimens must even be kept in separate enclosures, due to their solitary and aggressive nature. Territoriality is one of their most important characteristics.
The double-wattled cassowary (C. casuarius) is the most popular species in captivity, and it is fairly common in European and American zoos, where it is known for its unmistakable appearance. , only Weltvogelpark Walsrode in Germany has all three species of cassowaries in its collection: single-wattled cassowary (Casuarius unappendiculatus) and Bennett's cassowary (Casuarius bennetti). If subspecies are recognised, Weltvogelpark Walsrode has C. b. westermanni and C. u. rufotinctus.
Relationship with humans
Role in Papuan cultures and semi-domestication
There is evidence that the cassowary may have been domesticated by humans thousands of years before the chicken. Some New Guinea Highlands societies capture cassowary chicks and raise them as semi-tame poultry, for use in ceremonial gift exchanges and as food. They are the only indigenous Australasian animal known to have been partly domesticated by people prior to European arrival and colonization and by definition, the oldest form of domesticated animal and the largest domesticated bird. The Maring people of Kundagai sacrificed cassowaries (C. bennetti) in certain rituals. The Kalam people considered themselves related to cassowaries, and did not classify them as birds, but as kin. Consequently, they use the Pandanus register of the Kalam language when eating cassowary meat.
Studies on Pleistocene/early Holocene cassowary remains in Papua suggest that indigenous people at the time preferred to harvest eggs rather than adults. They seem to have regulated their consumption of these birds, possibly even collecting eggs and rearing young birds as one of the earliest forms of domestication.
Urbanisation of local cassowary population
In extremely urbanised areas where cassowaries used to naturally live such as in Queensland, Australia or in Port Moresby, Papua New Guinea, the local cassowary population had adapted to its less forested grounds. Increasing urbanisation has increased the likelihood of human-cassowary interaction, a potentially dangerous mix. Although cassowary populations have faced challenges in these urban areas in Northeastern Australia and parts of New Guinea, the cassowaries have proven to be surprisingly quite adaptable in contrast to the kiwis of New Zealand, potentially making them the largest urbanised birds in the world.
It was found that cassowaries in these urban environments changed their diets accordingly, with urbanised cassowaries actually consuming an even greater proportion of fruits from exotic plants (~30%) but still incorporating a significant proportion of fruits from native plants in their diet. Likewise, as aforementioned, the high concentration of human activity in the urban ecology also equates to a higher concentration of food diversity and food waste, with these 'urbanised' cassowaries foraging for food scraps, bird feeders and outdoor picnic/food venues without fear from humans or domesticated animals due to the birds' size and reputation.
Due to their omnivorous nature, cassowaries are able to eat all types of human food, including processed ones if interested, although fruit still remains their favourite choosings. A 2013 study from post-mortem investigations found that a combination of fruit scarcity and abundancy in human waste saw the diet of the cassowary intaking vast quantities of non-fruit items, this include fungus, carrion, meat, cheese, bones, pasta, chilli and tomato. The high concentration of human activity as well as vehicles, mixed with domesticated animals and less forest coverage, had also changed their behaviours. These 'city' cassowaries were shown to exist in a higher state of activity and rested less than individuals inhabiting more intact swathes of rainforest, actively moving between urban gardens and the rainforest. The study give evidence that these birds showed a surprising amount of flexible foraging strategy that has enabled them to persist in rainforest-fragmented landscapes.
Attacks
Cassowaries have a reputation for being dangerous to people and domestic animals. During World War II, American and Australian troops stationed in New Guinea were warned to steer clear of them. In his 1958 book Living Birds of the World, ornithologist Ernest Thomas Gilliard wrote:
The inner or second of the three toes is fitted with a long, straight, murderous nail which can sever an arm or eviscerate an abdomen with ease. There are many records of natives being killed by this bird.
This assessment of the danger posed by cassowaries has been repeated in print by authors, including Gregory S. Paul and Jared Diamond. A 2003 historical study of 221 cassowary attacks showed that 150 had been against humans; 75% of these had been from cassowaries that had been fed by people, 71% of the time the bird had chased or charged the victim, and 15% of the time they kicked. Of the attacks, 73% involved the birds expecting or snatching food, 5% involved defending their natural food sources, 15% involved defending themselves, and 7% involved defending their chicks or eggs. Only one human death was reported among those 150 attacks.
The first documented human death caused by a cassowary was on April 6, 1926. In Australia, 16-year-old Phillip McClean and his brother, age 13, came across a cassowary on their property and decided to try to kill it by striking it with clubs. The bird kicked the younger boy, who fell and ran away as his older brother struck the bird. The older McClean then tripped and fell to the ground. While he was on the ground, the cassowary kicked him in the neck, opening a wound that severed his jugular vein. The boy died of his injuries shortly thereafter.
Cassowary strikes to the abdomen are among the rarest of all, but in one case in 1995, a dog was kicked in the belly. The blow left no puncture, but severe bruising occurred. The dog later died from an apparent intestinal rupture.
Another human death due to a cassowary was recorded in Florida on April 12, 2019. The bird's owner, a 75-year-old man who had raised the animal, was apparently clawed to death after he fell to the ground.
| Biology and health sciences | Ratites | null |
59171 | https://en.wikipedia.org/wiki/Mumps | Mumps | Mumps is a highly contagious viral disease caused by the mumps virus. Initial symptoms of mumps are non-specific and include fever, headache, malaise, muscle pain, and loss of appetite. These symptoms are usually followed by painful swelling around the side of the face (the parotid glands, called parotitis), which is the most common symptom of a mumps infection. Symptoms typically occur 16 to 18 days after exposure to the virus. About one-third of people with a mumps infection do not have any symptoms (asymptomatic).
Complications are rare but include deafness and a wide range of inflammatory conditions, of which inflammation of the testes, breasts, ovaries, pancreas, meninges, and brain are the most common. Viral meningitis can occur in 1/4 of people with mumps. Testicular inflammation may result in reduced fertility and, rarely, sterility.
Humans are the only natural hosts of the mumps virus. The mumps virus is an RNA virus in the family Paramyxoviridae. The virus is primarily transmitted by respiratory secretions such as droplets and saliva, as well as via direct contact with an infected person. Mumps is highly contagious and spreads easily in densely populated settings. Transmission can occur from one week before the onset of symptoms to eight days after. During infection, the virus first infects the upper respiratory tract. From there, it spreads to the salivary glands and lymph nodes. Infection of the lymph nodes leads to the presence of the virus in the blood, which spreads the virus throughout the body. In places where mumps is common, it can be diagnosed based on clinical presentation. In places where mumps is less common, however, laboratory diagnosis using antibody testing, viral cultures, or real-time reverse transcription polymerase chain reaction may be needed.
There is no specific treatment for mumps, so treatment is supportive and includes rest and pain relief. Mumps infection is usually self-limiting, coming to an end as the immune system clears the infection. Infection can be prevented with vaccination. The MMR vaccine is a safe and effective vaccine to prevent mumps infections and is used widely around the world. The MMR vaccine also protects against measles and rubella. The spread of the disease can also be prevented by isolating infected individuals.
Mumps historically has been a highly prevalent disease, commonly occurring in outbreaks in densely crowded spaces. In the absence of vaccination, infection normally occurs in childhood, most frequently at the ages of 5–9. Symptoms and complications are more common in males and more severe in adolescents and adults. Infection is most common in winter and spring in temperate climates, whereas no seasonality is observed in tropical regions. Written accounts of mumps have existed since ancient times, and the cause of mumps, the mumps virus, was discovered in 1934. By the 1970s, vaccines had been created to protect against infection, and countries that have adopted mumps vaccination have seen a near-elimination of the disease. In the 21st century, however, there has been a resurgence in the number of cases in many countries that vaccinate, primarily among adolescents and young adults, due to multiple factors such as waning vaccine immunity and opposition to vaccination.
Etymology
The word "mumps" was first attested circa 1600 and is the plural form of "mump", meaning "grimace", originally a verb meaning "to whine or mutter like a beggar". The disease was likely called mumps due to the swelling caused by mumps parotitis, reflecting its impact on facial expressions and the painful, difficult swallowing that it causes. "Mumps" was also used starting from the 17th century to mean "a fit of melancholy, sullenness, silent displeasure". Mumps is sometimes called "epidemic parotitis".
History
According to Chinese medical literature, mumps was recorded as far back as 640 B.C. The Greek physician Hippocrates documented an outbreak on the island of Thasos in approximately 410 B.C. and provided a fuller description of the disease in the first book of Epidemics in the Corpus Hippocraticum. In modern times, the disease was first described scientifically in 1790 by British physician Robert Hamilton in Transactions of the Royal Society of Edinburgh. During the First World War, mumps was one of the most debilitating diseases among soldiers. In 1934, the etiology of the disease, the mumps virus, was discovered by Claude D. Johnson and Ernest William Goodpasture. They found that rhesus macaques exposed to saliva taken from humans in the early stages of the disease developed mumps. Furthermore, they showed that mumps could then be transferred to children via filtered and sterilized, bacteria-less preparations of macerated monkey parotid tissue, showing that it was a viral disease.
In 1945, the mumps virus was isolated for the first time. Just a few years later, in 1948, an inactivated vaccine using killed viruses was invented. This vaccine provided only short-term immunity and was later discontinued. It was replaced in the 1970s with vaccines that have live but weakened viruses, which are more effective at providing long-term immunity than the inactivated vaccine. The first of these vaccines was Mumpsvax, licensed on 30 March 1967, which used the Jeryl Lynn strain. Maurice Hilleman created this vaccine using the strain taken from his five-year-old daughter, Jeryl Lynn. Mumpsvax was recommended for use in 1977, and the Jeryl Lynn strain continues to be used.
Hilleman worked to combine the attenuated mumps vaccines with the measles and rubella vaccines, creating the MMR-1 vaccine. In 1971, a newer version, MMR-2, was approved for use by the US Food and Drug Administration. In the 1980s, the benefit of multiple doses was recognized, so a two-dose immunization schedule was widely adopted. With MMR-2, four other MMR vaccines have been created since the 1960s: Triviraten, Morupar, Priorix, and Trimovax. Since the mid-2000s, two MMRV vaccines have been in use: Priorix-Tetra and ProQuad.
The United States began to vaccinate against mumps in the 1960s, with other countries following suit. From 1977 to 1985, 290 cases per 100,000 people were diagnosed each year worldwide. Although few countries recorded mumps cases after they began vaccination, those that did reported dramatic declines. From 1968 to 1982, cases declined by 97% in the U.S., and in Finland cases were reduced to less than one per 100,000 people per year, and a decline from 160 cases per 100,000 to 17 per 100,000 per year in England was observed from 1989 to 1995. By 2001, there had been a 99.9% reduction in the number of cases in the U.S. and similar near-elimination in other vaccinating countries.
In Japan in 1993, concerns over the rates of aseptic meningitis following MMR vaccination with the Urabe strain prompted the removal of MMR vaccines from the national immunization program, resulting in a dramatic increase in the number of cases. Japan provides voluntary mumps vaccination separately from measles and rubella. Starting in the mid-1990s, controversies surrounding the MMR vaccine emerged. One paper connected the MMR vaccine to Crohn's disease in 1995, and another in 1998 connected it to autism spectrum disorders and inflammatory bowel disease. These papers are now considered to be fraudulent and incorrect, and no association between the MMR vaccine and the aforementioned conditions has been identified. Despite this, their publication led to a significant decline in vaccination rates, ultimately causing measles, mumps, and rubella to reemerge in places with lowered vaccination rates.
Outbreaks in the 21st century include more than 300,000 cases in China in 2013 and more than 56,000 cases in England and Wales in 2004–2005. In the latter outbreak, most cases were reported in 15–24 year olds who were attending colleges and universities. This age group was thought to be vulnerable to infection because of the MMR vaccine controversies when they should have been vaccinated or MMR vaccine shortages that had also occurred at that time. Similar outbreaks in densely crowded environments have frequently occurred in many other countries, including the U.S., the Netherlands, Sweden, and Belgium.
Resurgence
In the 21st century, mumps has reemerged in many places that vaccinate against it, causing recurrent outbreaks. These outbreaks have largely affected adolescents and young adults in densely crowded spaces, such as schools, sports teams, religious gatherings, and the military, and it is expected that outbreaks will continue to occur. The cause of this reemergence is subject to debate, and various factors have been proposed, including waning immunity from vaccination, low vaccination rates, vaccine failure, and potential antigenic variation of the mumps virus.
Waning immunity from vaccines is likely the primary cause of the mumps resurgence. In the past, subclinical natural infections provided boosts to immunity similar to vaccines. As time went on with vaccine use, these asymptomatic infections declined in frequency, likely leading to a reduction in long-term immunity against mumps. With less long-term immunity, the effects of waning vaccine immunity became more prominent, and vaccinated individuals have frequently fallen ill from mumps. A third dose of the vaccine provided in adolescence has been considered to address this as some studies support this. Other research indicates that a third dose may be useful only for short-term immunity in responding to outbreaks, which is recommended for at-risk persons by the Advisory Committee on Immunization Practices of the Centers for Disease Control and Prevention.
Low vaccination rates have been implicated as the cause of some outbreaks in the UK, Canada, Sweden, and Japan, whereas outbreaks in other places, such as the U.S., the Czech Republic, and the Netherlands, have occurred mainly among the vaccinated. Compared to the measles and rubella vaccines, mumps vaccines appear to have a relatively high failure rate, varying depending on the vaccine strain. This has been addressed by providing two vaccine doses, supported by recent outbreaks among the vaccinated having primarily occurred among those who received only one dose. Lastly, certain mumps virus lineages are highly divergent genetically from vaccine strains, which may cause a mismatch between protection against vaccine strains and non-vaccine strains, though research is inconclusive on this matter.
Signs and symptoms
Common symptoms
The incubation period, the time between the start of an infection and when symptoms begin to show, is about 7–25 days, averaging 16–18 days. 20–40% of infections are asymptomatic or are restricted to mild respiratory symptoms, sometimes with a fever. Over the course of the disease, three distinct phases are recognized: prodromal, early acute, and established acute. The prodromal phase typically has non-specific, mild symptoms such as a low-grade fever, headache, malaise, muscle pain, loss of appetite, and sore throat. In the early acute phase, as the mumps virus spreads throughout the body, systemic symptoms emerge. Most commonly, parotitis occurs during this time period. During the established acute phase, orchitis, meningitis, and encephalitis may occur, and these conditions are responsible for the bulk of mumps morbidity.
The parotid glands are salivary glands situated on the sides of the mouth in front of the ears. Inflammation of them, called parotitis, is the most common mumps symptom and occurs in about 90% of symptomatic cases and 60–70% of total infections. During mumps parotitis, usually both the left and right parotid glands experience painful swelling, with unilateral swelling in a small percentage of cases. Parotitis occurs 2–3 weeks after exposure to the virus, within two days of developing symptoms, and usually lasts 2–3 days, but it may last as long as a week or longer.
In 90% of parotitis cases, swelling on one side is delayed rather than both sides swelling in unison. The parotid duct, which is the opening that provides saliva from the parotid glands to the mouth, may become red, swollen, and filled with fluid. Parotitis is usually preceded by local tenderness and occasionally earache. Other salivary glands, namely the submandibular, and sublingual glands, may also swell. Inflammation of these glands is rarely the only symptom.
Complications
Outside of the salivary glands, inflammation of the testes, called orchitis, is the most common symptom of infection. Pain, swelling, and warmness of a testis appear usually 1–2 weeks after the onset of parotitis but can occur up to six weeks later. During mumps orchitis, the scrotum is tender and inflamed. It occurs in 10–40% of pubertal and post-pubertal males who contract mumps. Usually, mumps orchitis affects only one testis but in 10–30% of cases both are affected. Mumps orchitis is accompanied by inflammation of the epididymis, called epididymitis, about 85% of the time, typically occurring before orchitis. The onset of mumps orchitis is associated with a high-grade fever, vomiting, headache, and malaise. In prepubertal males, orchitis is rare as symptoms are usually restricted to parotitis.
A variety of other inflammatory conditions may also occur as a result of mumps virus infection, including:
Mastitis, inflammation of the breasts, in up to about 30% of post-pubertal women
Oophoritis, inflammation of an ovary, in 5–10% of post-pubertal women, which usually presents as pelvic pain
Aseptic meningitis, inflammation of the meninges, in 5–10% of cases and 4–6% of those with parotitis, typically occurring 4–10 days after the onset of symptoms. Mumps meningitis can also occur up to one week before parotitis as well as in the absence of parotitis. It is commonly accompanied by fever, headache, vomiting, and neck stiffness.
Pancreatitis, inflammation of the pancreas, in about 4% of cases, which causes severe pain and tenderness in the upper abdomen below the ribs
Encephalitis, inflammation of the brain, in less than 0.5% of cases. People who experience mumps encephalitis typically experience a fever, altered consciousness, seizures, and weakness. Like meningitis, mumps encephalitis can occur in the absence of parotitis.
Meningoencephalitis, inflammation of the brain and its surrounding membranes. Mumps meningoencephalitis is commonly accompanied by fever 97% of the time, vomiting 94% of the time, and headache 88.8% of the time.
Nephritis, inflammation of the kidneys, which is rare because kidney involvement in mumps is usually benign but leads to presence of the virus in urine
Inflammation of the joints (arthritis), which may affect at least five joints (polyarthritis), multiple nerves in the peripheral nervous system (polyneuritis), pneumonia, gallbladder without gallstones (acalculous cholecystitis), cornea and uveal tract (keratouveitis), thyroids (thyroiditis), liver (hepatitis), retina (retinitis), and corneal endothelium (corneal endothelitis), all of which are rare
Recurrent sialadenitis, inflammation of the salivary glands, which is frequent
A relatively common complication is deafness, which occurs in about 4% of cases. Mumps deafness is often accompanied by vestibular symptoms such as vertigo and repetitive, uncontrolled eye movements. Based on electrocardiographic abnormalities in the infected, MuV also likely infects cardiac tissue, but this is usually asymptomatic. Rarely, myocarditis and pericarditis can occur. Fluid buildup in the brain, called hydrocephalus, has also been observed. In the first trimester of pregnancy, mumps may increase the risk of miscarriage. Otherwise, mumps is not associated with birth defects.
Other rare complications of infection include: paralysis, seizures, cranial nerve palsies, cerebellar ataxia, transverse myelitis, ascending polyradiculitis, a polio-like disease, arthropathy, autoimmune hemolytic anemia, idiopathic thrombocytopenic purpura, Guillain–Barré syndrome, post-infectious encephalitis encephalomyelitis, and hemophagocytic syndrome. At least one complication occurs in combination with the standard mumps symptoms in up to 42% of cases. Mumps has also been connected to the onset of type 1 diabetes, and, relatedly, the mumps virus is able to infect and replicate in insulin-producing beta cells. Among children, seizures occur in about 20–30% of cases involving the central nervous system.
Cause
Mumps is caused by the mumps virus (MuV), scientific name Mumps orthorubulavirus, which belongs to the Orthorubulavirus genus in the Paramyxoviridae family of viruses. Humans are the only natural host of the mumps virus. MuV's genome is made of RNA and contains seven genes that encode nine proteins. In MuV particles, the genome is encased by a helical capsid. The capsid is surrounded by a viral envelope that has spikes protruding from its surface. MuV particles are pleomorphic in shape and range from 100 to 600 nanometers in diameter.
The replication cycle of MuV begins when the spikes on its surface bond to a cell, which then causes the envelope to fuse with the host cell's cell membrane, releasing the capsid into the host cell's cytoplasm. Upon entry, the viral RNA-dependent RNA polymerase (RdRp) transcribes messenger RNA (mRNA) from the genome, which is then translated by the host cell's ribosomes to synthesize viral proteins. RdRp then begins replicating the viral genome to produce progeny. Viral spike proteins fuse into the host cell's membrane, and new virions are formed at the sites beneath the spikes. MuV then utilizes host cell proteins to leave the host cell by budding from its surface, using the host cell's membrane as the viral envelope.
Twelve genotypes of MuV are recognized, named genotypes A to N, excluding E and M. These genotypes vary in frequency from region to region. For example, genotypes C, D, H, and J are more common in the western hemisphere, whereas genotypes F, G, and I are more common in Asia, although genotype G is considered to be a global genotype. Genotypes A and B have not been observed in the wild since the 1990s. MuV has just one serotype, so antibodies to one genotype are functional against all genotypes. MuV is a relatively stable virus and is unlikely to experience antigenic shifting that may cause new strains to emerge.
Transmission
The mumps virus is mainly transmitted by inhalation or oral contact with respiratory droplets or secretions. In experiments, mumps could develop after inoculation either via the mouth or the nose. Respiratory transmission is also supported by the presence of MuV in cases of respiratory illness without parotitis, detection in nasal samples, and transmission between people in close contact. MuV is excreted in saliva from approximately one week before to eight days after the onset of symptoms, peaking at the onset of parotitis, though it has also been identified in the saliva of asymptomatic individuals.
Mother-to-child transmission has been observed in various forms. In non-human primates, placental transmission has been observed, which is supported by the isolation of MuV from spontaneous and planned aborted fetuses during maternal mumps. MuV has also been isolated from newborns whose mother was infected. While MuV has been detected in breast milk, it is unclear if the virus can be transmitted through it. Other manners of transmission include direct contact with infected droplets or saliva, fomites contaminated by saliva and possibly urine. Most transmissions likely occur before the development of symptoms and up to five days after such time.
In susceptible populations, a single case can cause up to twelve new ones. The period when a person is contagious lasts from two days before the onset of symptoms to nine days after symptoms have ceased. Asymptomatic carriers of the mump virus can also transmit the virus. These factors are thought to be reasons why controlling the spread of mumps is difficult. Furthermore, reinfection can occur after natural infection or vaccination, indicating that lifelong immunity is not guaranteed after infection. Vaccinated individuals who are infected appear to be less contagious than the unvaccinated.
The average number of new cases generated from a single case in a susceptible population called the basic reproduction number, is 4–7. Given this, it is estimated that a vaccination rate between 79 and 100% is needed to achieve herd immunity. Outbreaks continue to occur in places that have vaccination rates exceeding 90%, however, suggesting that other factors may influence disease transmission. Outbreaks that have occurred in these vaccinated communities typically occur in highly crowded areas such as schools and military dormitories.
Pathogenesis
Many aspects of the pathogenesis of mumps are poorly understood and are inferred from clinical observations and experimental infections in laboratory animals. These animal studies may be unreliable due to unnatural methods of inoculation. Following exposure, the virus infects epithelial cells in the upper respiratory tract that express sialic acid receptors on their surface. After infection, the virus spreads to the parotid glands, causing the signature parotitis. It is thought that shortly after infection the virus spreads to lymph nodes, in particular T-cells and viruses in the blood, called viremia. Viremia lasts for 7–10 days, during which MuV spreads throughout the body.
In mumps orchitis, infection leads to: parenchymal edema; congestion, or separation, of the seminiferous tubules; and perivascular infiltration by lymphocytes. The tunica albuginea forms a barrier against edema, causing an increase in intratesticular pressure that causes necrosis of the seminiferous tubules. The seminiferous tubules also experience hyalinization, i.e. degeneration into a translucent glass-like substance, which can cause fibrosis and atrophy of the testes.
In up to half of cases, MuV infiltrates the central nervous system (CNS), where it may cause meningitis, encephalitis, or hydrocephalus. Mumps is rarely fatal, so few post-mortem analyses have been done to analyze CNS involvement. Of these, fluid buildup, congestion, and hemorrhaging in the brain, white blood cell infiltration in the perivascular spaces in the brain, reactive changes to glial cells and damage to the myelin sheaths surrounding neurons were observed. Neurons appear to be relatively unaffected.
In laboratory tests on rodents, MuV appears to enter the CNS first through cerebrospinal fluid (CSF), then spreading to the ventricular system. There, MuV replicates in ependymal cells that line the ventricles, which allows the virus to enter the brain parenchyma. This often leads to MuV infecting pyramidal cells in the cerebral cortex and hippocampus. Infected ependymal cells become inflamed, lose their cilia, and collapse into CSF, which may be the cause of the narrowing of the cerebral aqueduct thought to cause mumps hydrocephalus.
In humans, mumps hydrocephalus may be due to obstruction of the cerebral aqueduct with dilatation of the lateral and third ventricles, obstruction of the interventricular foramina, or obstruction of the median and lateral apertures. Ependymal cells have been isolated from CSF of mumps patients, suggesting that animals and humans share hydrocephalus pathogenesis. Hydrocephalus has also been observed in the absence of canal obstruction, however, indicating that obstruction may be a result of external compression by edematous tissue and not related to hydrocephalus.
Deafness from mumps may be caused by MuV infection in CSF, which has contact with the perilymph of the inner ear, possibly leading to infection of the cochlea, or it may occur as a result of inner ear infection via viremia that leads to inflammation in the endolymph. Hearing loss may also be caused indirectly by the immune response. In animal studies, MuV has been isolated from the vestibular ganglion, which may explain vestibular symptoms such as vertigo that often co-occur with deafness.
Immune response
Even though MuV has just one serotype, significant variation in the quantity of genotype-specific sera needed to neutralize different genotypes in vitro has been observed. Neutralizing antibodies in the salivary glands may be important in restricting MuV replication and transmission via saliva, as the level of viral secretion in saliva inversely correlates to the quantity of MuV-specific IgA produced. The neutralizing ability of salivary IgA appears to be greater than serum IgG and IgM.
It has been proposed that symptomatic infections in the vaccinated may be because memory T lymphocytes generated as a result of vaccination may be necessary but insufficient for protection. The immune system in general appears to have a relatively weak response to the mumps virus, indicated by various measures: antibody production appears to be predominately directed toward non-neutralizing viral proteins, and there may be a low quantity of MuV-specific memory B lymphocytes. The amount of antibodies needed to confer immunity is unknown.
Diagnosis
In places where mumps is widespread, diagnosis can be made based on the development of parotitis and a history of exposure to someone with mumps. In places where mumps is less common because parotitis has other causes, laboratory diagnosis may be needed to verify mumps infection. A differential diagnosis may be used to compare symptoms to other diseases, including allergic reaction, mastoiditis, measles, and pediatric HIV infection and rubella. MuV can be isolated from saliva, blood, the nasopharynx, salivary ducts, and seminal fluid within one week of the onset of symptoms, as well as from cell cultures. In meningitis cases, MuV can be isolated from CSF. In CNS cases, a lumbar puncture may be used to rule out other potential causes, which shows normal opening pressure, more than ten leukocytes per cubic millimeter, elevated lymphocyte count in CSF, polymorphonuclear leukocytes up to 25% of the time, often a mildly elevated protein level, and a slightly reduced CSF glucose to blood glucose ratio up to 30% of the time.
Mumps-specific IgM antibodies in serum or oral fluid specimens can be used to identify mumps. IgM quantities peak up to eight days after the onset of symptoms, and IgM can be measured by enzyme-linked immunosorbent assays (ELISA) 7–10 days after the onset of symptoms. Sensitivity to IgM testing is variable, ranging from as low as 24–51% to 75% in the first week and 100% thereafter. Throughout infection, IgM titers increase four-fold between the acute phase and recovery. False negatives can occur in people previously infected or vaccinated, in which case a rise of serum IgG may be more useful for diagnosis. False positives can occur after infection of parainfluenza viruses1 and 3 and Newcastle disease virus as well as recently after mumps vaccination.
Antibody titers can also be measured with complement fixation tests, hemagglutination assays, and neutralization tests. In vaccinated people, antibody-based diagnosis can be difficult since IgM oftentimes cannot be detected in acute-phase serum samples. In these instances, it is easier to identify MuV RNA from oral fluid, a throat swab, or urine. In meningitis cases, MuV-specific IgM can be found in CSF in half of cases, and IgG in 30–90%, sometimes lasting for more than a year with increased white blood cell count. These findings are not associated with an increased risk of long-term complications. Most parotitis cases have elevated white blood cell count in CSF.
Real-time reverse transcription polymerase chain reaction (rRT-PCR) can be used to detect MuV RNA from the first day symptoms appear, declining over the next 8–10 days. rRT-PCR of saliva is typically positive from 2–3 days before parotitis develops to 4–5 days after and has a sensitivity of about 70%. Since MuV replicates in kidneys, viral culture and RNA detection in urine can be used for diagnosis up to two weeks after symptoms begin, though rRT-PCR used to identify the virus in urine has a very low sensitivity compared to virus cultures at below 30%. In meningoencephalitis cases, a nested RT-PCR can detect MuV RNA in CSF up to two years after infection.
In sialadenitis cases, imaging shows an enlargement of the salivary glands, fat stranding, and thickening of the superficial cervical fascia and platysma muscles, which are situated on the front side of the neck. If parotitis occurs only on one side, then detection of mumps-specific IgM antibodies, IgG titer, or PCR is required for diagnosis. In cases of pancreatitis, there may be elevated levels of lipase or amylase, an enzyme found in saliva and the pancreas.
Mumps orchitis is usually diagnosed by white blood cell count, with normal differential white blood cell counts. A complete blood count can show above or below-average white blood cell count and an elevated C-reactive protein level. Urine analysis can exclude bacterial infections. If orchitis is present with normal urine analysis, negative urethral cultures, and negative midstream urine, then that can indicate mumps orchitis. Ultrasounds typically show diffuse hyper-vascularity, increased volume of the testes and epididymis, lower than usual ability to return ultrasound signals, swelling of the epididymis, and formation of hydroceles. Echo color Doppler ultrasound is more effective at detecting orchitis than ultrasound alone.
Prevention
Mumps is preventable with vaccination. Mumps vaccines use live attenuated viruses. Most countries include mumps vaccination in their immunization programs, and the MMR vaccine, which also protects against measles and rubella, is the most commonly used mumps vaccine. Mumps vaccination can also be done on its own and as a part of the MMRV vaccine, which also protects against measles, rubella, chickenpox, and shingles. More than 120 countries have adopted mumps vaccination, but coverage remains low in most African, South Asian, and Southeast Asian countries. In countries that have implemented mumps vaccination, significant declines in mumps cases and complications caused by infection such as encephalitis have been observed. Mumps vaccines are typically administered in early childhood, but may also be given in adolescence and adulthood if need be. Vaccination is expected to be capable of neutralizing wild-type MuVs, which are not included in the vaccine, since they do not appear to evade vaccine-derived immunity.
A variety of virus strains have been used in mumps vaccines, including the Jeryl Lynn (JL), Leningrad-3, Leningrad-3-Zagreb (L-Zagreb), Rubini, and Urabe AM9 strains. Some other less prominent strains exist that are typically confined to individual countries. These include the Hoshino, Miyahara, Torii, and NK M-46 strains that have been produced in Japan and the S-12 strain, which is used by Iran. Mild adverse reactions are relatively common, including fever and rash, but aseptic meningitis also occurs at varying rates. Other rare adverse reactions include meningoencephalitis, parotitis, deafness from inner ear damage, orchitis, and pancreatitis. Safety and effectiveness vary by vaccine strain:
Rubini is safe but because of its low effectiveness in outbreaks, its use has been abandoned.
JL is relatively safe and has a relatively high effectiveness. However, the effectiveness is significantly lower in outbreaks. A modified version of JL vaccines is RIT 4385, which is also considered safe.
Urabe and Leningrad-3 are both at least as effective as JL, but are less safe.
L-Zagreb, a modified version of Leningrad-3, is considered safe and effective, including in outbreaks.
Mumps protection from the MMR vaccine is higher after two doses than one and is estimated to be between 79% and 95%, lower than the degree of protection against measles and rubella. This, however, has still been sufficient to nearly eliminate mumps in countries that vaccinate against it as well as significantly reduce frequencies of complications among the vaccinated. If at least one dose is received, then hospitalization rates are reduced by an estimated 50% among the infected. Compared to the MMR vaccine, the MMRV vaccine appears to be less effective in terms of providing mumps protection. A difficulty in assessing vaccine effectiveness is that there is no clear correlate of immunity, so it is not possible to predict if a person has acquired immunity from the vaccine.
There is a lack of data on the effectiveness of a third dose of the MMR vaccine. In an outbreak in which a third dose was administered, it was unclear if it had any effect on reducing disease incidence, and it only appeared to boost antibodies in those who previously had little or no antibodies to mumps. Contraindications for mumps vaccines include prior allergic reaction to any ingredients or to neomycin, pregnancy, immunosuppression, a moderate or severe illness, having received a blood product recently, and, for MMRV vaccines specifically, a personal or familial history of seizures. It is also advised that women not become pregnant in the four weeks after MMR vaccination. No effective prophylaxis exists for mumps after one has been exposed to the virus, so vaccination or receiving immunoglobulin after exposure does not prevent progression to illness.
For people who are infected or suspected to be infected, isolation is important in preventing the spread of the disease. This includes abstaining from school, childcare, work, and other settings in which people gather together. In healthcare settings, it is recommended that healthcare workers use precautions such as face masks to reduce the likelihood of infection and to abstain from work if they develop mumps. Additional measures taken in health care facilities include reducing wait times for mumps patients, having mumps patients wear masks, and cleaning and disinfecting areas that mumps patients use. The virus can be inactivated using formalin, ether, chloroform, heat, or ultraviolet light.
Treatment
Mumps is usually self-limiting, and no specific antiviral treatments exist for it, so treatment is aimed at alleviating symptoms and preventing complications. Non-medicinal ways to manage the disease include bed rest, using ice or heat packs on the neck and scrotum, consuming more fluids, eating soft food, and gargling with warm salt water. Anti-fever medications may be used during the febrile period, excluding aspirin when given to children, which may cause Reye syndrome. Analgesics may also be provided to control pain from mumps inflammatory conditions. For seizures, anticonvulsants may be used. In severe neurological cases, ventilators may be used to support breathing.
Intramuscular mumps immunoglobulin may be of benefit when administered early in some cases, but it has not shown benefit in outbreaks. Although not recommended, intravenous immunoglobulin therapy may reduce the rates of some complications. Antibiotics may be used as a precaution in cases in which bacterial infection cannot be ruled out as well as to prevent secondary bacterial infection. Autoimmune-based disorders connected to mumps are treatable with intravenous immunoglobulin.
Various types of treatment for mumps orchitis have been used, but no specific treatment is recommended due to each method's limitations. These measures are primarily based on relieving testicular pain and reducing intratesticular pressure to reduce the likelihood of testicular atrophy. Interferon-α2α interferes with viral replication, so it has been postulated to be useful in preventing testicular damage and infertility. Interferon alfa-2b may reduce the duration of symptoms and incidence of complications. In cases of hydrocele formation, excess fluid can be removed.
Acupuncture has been used fairly widely in China to treat children who have mumps, however, no high-quality trials have been conducted to determine the safety or effectiveness of this treatment approach.
Prognosis
The prognosis for most people who experience mumps is excellent as long-term complications and death are rare. Hospitalization is typically not required. Mumps is usually self-limiting and symptoms resolve spontaneously within two weeks as the immune system clears the virus from the body. In high-risk groups such as immunocompromised persons, the prognosis is considered to be the same as for other groups. For most people, infection leads to lifelong immunity against future infection. Reinfections appear to be more mild and atypical than the first infection. The overall case-fatality rate of mumps is 1.6–3.8 people per 10,000, and these deaths typically occur in those who develop encephalitis.
Mumps orchitis typically resolves within two weeks. In 20% of cases, the testicles may be tender for a few more weeks. Atrophy, or reduction of size, of the involved testicle occurs in 30–50% of orchitis cases, which may lead to abnormalities in sperm creation and fertility such as low sperm count, absence of sperm in semen, reduced sperm motility, reduced fertility (hypofertility) in 13% of cases, and rarely sterility. Hypofertility can, however, occur in cases without atrophy. Abnormalities in sperm creation can persist for months to years after recovery from the initial infection, the length of which increases as the severity of orchitis increases. Examination of these cases shows decreased testicular volume, tenderness of the testicles, and a feeling of inconsistency when handling the testicles. Infertility is linked to severe cases of orchitis affecting both testes followed by testicular atrophy, which may develop up to one year after the initial infection. Of bilateral orchitis cases, 30–87% experience infertility. There is a weak association between orchitis and later development of epididymitis and testicular tumors.
Mumps meningitis typically resolves within 3–10 days without long-term complications. In meningoencephalitis cases, higher protein levels in CSF and a lower CSF glucose to blood glucose ratio are associated with longer periods of hospitalization. Approximately 1% of those whose CNS is affected die from mumps. Post-infectious encephalitis tends to be relatively mild, whereas post-infectious encephalomyelitis has a case-fatality rate of up to ten percent. Most cases of mumps deafness affect just one ear and are temporary, but permanent hearing loss occurs in 0.005% of infections. Myocarditis and pericarditis that occur as a result of mumps may lead to endocardial fibroelastosis, i.e. thickening of the endocardium. With extreme rarity, infertility and premature menopause have occurred as a result of mumps oophoritis.
Epidemiology
Clinical age and immunity
Mumps is found worldwide. In the absence of vaccination against mumps there are between 100 and 1,000 cases per 100,000 people each year, i.e. 0.1% to 1.0% of the population are infected each year. The number of cases peaks every 2–5 years, with incidence highest in children 5–9 years old. According to seroconversion surveys done before the start of mumps vaccination, a sharp increase in mumps antibody levels at age 2–3 was observed.
Furthermore, 50% of 4–6 year olds, 90% of 14–15 year olds, and 95% of adults had tested positive to prior exposure to mumps, indicating that nearly all people are eventually infected in unvaccinated populations.
Prior to the start of vaccination, mumps accounted for ten percent of meningitis cases and about a third of encephalitis cases. Worldwide, mumps is the most common cause of inflammation of the salivary glands. In children, mumps is the most common cause of deafness in one ear in cases when the inner ear is damaged. Asymptomatic infections are more common in adults, and the rate of asymptomatic infections is very high, up to two-thirds, in vaccinated populations. Mumps vaccination has the effect of increasing the average age of the infected in vaccinated populations that have not previously experienced a mumps outbreak. While infection rates appear to be the same in males and females, males appear to experience symptoms and complications, including neurological involvement, at a higher rate than females. Symptoms are more severe in adolescents and adults than in children.
Settings of outbreaks
It is common for outbreaks of mumps to occur. These outbreaks typically occur in crowded spaces where the virus can spread from person to person easily, such as schools, military barracks, prisons, and sports clubs. Since the introduction of vaccines, the frequency of mumps has declined dramatically, as have complications caused by mumps. The epidemiology in countries that vaccinate reflects the number of doses administered, age at vaccination, and vaccination rates. If vaccine coverage is insufficient, then herd immunity may be unobtainable and the average age of infection will increase, leading to an increase in the prevalence of complications. Risk factors include age, exposure to a person with mumps, compromised immunity, time of year, travel history, and vaccination status. Mumps vaccination is less common in developing countries, which consequently have higher rates of mumps.
Cases peak in different seasons of the year in different regions. In temperate climates, cases peak in winter and spring, whereas in tropical regions no seasonality is observed. Additional research has shown that mumps increases in frequency as temperature and humidity increase. The seasonality of mumps is thought to be caused by several factors: fluctuation in the human immune response due to seasonal factors, such as changes in melatonin levels; behavior and lifestyle changes, such as school attendance and indoor crowding; and meteorological factors such as changes in temperature, brightness, wind, and humidity.
| Biology and health sciences | Infectious disease | null |
59217 | https://en.wikipedia.org/wiki/Quadratic%20formula | Quadratic formula | In elementary algebra, the quadratic formula is a closed-form expression describing the solutions of a quadratic equation. Other ways of solving quadratic equations, such as completing the square, yield the same solutions.
Given a general quadratic equation of the form , with representing an unknown, and coefficients , , and representing known real or complex numbers with , the values of satisfying the equation, called the roots or zeros, can be found using the quadratic formula,
where the plus–minus symbol "" indicates that the equation has two roots. Written separately, these are:
The quantity is known as the discriminant of the quadratic equation. If the coefficients , , and are real numbers then when , the equation has two distinct real roots; when , the equation has one repeated real root; and when , the equation has no real roots but has two distinct complex roots, which are complex conjugates of each other.
Geometrically, the roots represent the values at which the graph of the quadratic function , a parabola, crosses the -axis: the graph's -intercepts. The quadratic formula can also be used to identify the parabola's axis of symmetry.
Derivation by completing the square
The standard way to derive the quadratic formula is to apply the method of completing the square to the generic quadratic equation . The idea is to manipulate the equation into the form for some expressions and written in terms of the coefficients; take the square root of both sides; and then isolate .
We start by dividing the equation by the quadratic coefficient , which is allowed because is non-zero. Afterwards, we subtract the constant term to isolate it on the right-hand side:
The left-hand side is now of the form , and we can "complete the square" by adding a constant to obtain a squared binomial . In this example we add to both sides so that the left-hand side can be factored (see the figure):
Because the left-hand side is now a perfect square, we can easily take the square root of both sides:
Finally, subtracting from both sides to isolate produces the quadratic formula:
Equivalent formulations
The quadratic formula can equivalently be written using various alternative expressions, for instance
which can be derived by first dividing a quadratic equation by , resulting in , then substituting the new coefficients into the standard quadratic formula. Because this variant allows re-use of the intermediately calculated quantity , it can slightly reduce the arithmetic involved.
Square root in the denominator
A lesser known quadratic formula, first mentioned by Giulio Fagnano, describes the same roots via an equation with the square root in the denominator (assuming ):
Here the minus–plus symbol "" indicates that the two roots of the quadratic equation, in the same order as the standard quadratic formula, are
This variant has been jokingly called the "citardauq" formula ("quadratic" spelled backwards).
When has the opposite sign as either or , subtraction can cause catastrophic cancellation, resulting in poor accuracy in numerical calculations; choosing between the version of the quadratic formula with the square root in the numerator or denominator depending on the sign of can avoid this problem. See below.
This version of the quadratic formula is used in Muller's method for finding the roots of general functions. It can be derived from the standard formula from the identity , one of Vieta's formulas. Alternately, it can be derived by dividing each side of the equation by to get , applying the standard formula to find the two roots , and then taking the reciprocal to find the roots of the original equation.
Other derivations
Any generic method or algorithm for solving quadratic equations can be applied to an equation with symbolic coefficients and used to derive some closed-form expression equivalent to the quadratic formula. Alternative methods are sometimes simpler than completing the square, and may offer interesting insight into other areas of mathematics.
Completing the square by Śrīdhara's method
Instead of dividing by to isolate , it can be slightly simpler to multiply by instead to produce , which allows us to complete the square without need for fractions. Then the steps of the derivation are:
Multiply each side by .
Add to both sides to complete the square.
Take the square root of both sides.
Isolate .
Applying this method to a generic quadratic equation with symbolic coefficients yields the quadratic formula:
This method for completing the square is ancient and was known to the 8th–9th century Indian mathematician Śrīdhara. Compared with the modern standard method for completing the square, this alternate method avoids fractions until the last step and hence does not require a rearrangement after step 3 to obtain a common denominator in the right side.
By substitution
Another derivation uses a change of variables to eliminate the linear term. Then the equation takes the form in terms of a new variable and some constant expression , whose roots are then .
By substituting into , expanding the products and combining like terms, and then solving for , we have:
Finally, after taking a square root of both sides and substituting the resulting expression for back into the familiar quadratic formula emerges:
By using algebraic identities
The following method was used by many historical mathematicians:
Let the roots of the quadratic equation be and . The derivation starts from an identity for the square of a difference (valid for any two complex numbers), of which we can take the square root on both sides:
Since the coefficient , we can divide the quadratic equation by to obtain a monic polynomial with the same roots. Namely,
This implies that the sum and the product . Thus the identity can be rewritten:
Therefore,
The two possibilities for each of and are the same two roots in opposite order, so we can combine them into the standard quadratic equation:
By Lagrange resolvents
An alternative way of deriving the quadratic formula is via the method of Lagrange resolvents, which is an early part of Galois theory.
This method can be generalized to give the roots of cubic polynomials and quartic polynomials, and leads to Galois theory, which allows one to understand the solution of algebraic equations of any degree in terms of the symmetry group of their roots, the Galois group.
This approach focuses on the roots themselves rather than algebraically rearranging the original equation. Given a monic quadratic polynomial assume that and are the two roots. So the polynomial factors as
which implies and .
Since multiplication and addition are both commutative, exchanging the roots and will not change the coefficients and : one can say that and are symmetric polynomials in and . Specifically, they are the elementary symmetric polynomials – any symmetric polynomial in and can be expressed in terms of and instead.
The Galois theory approach to analyzing and solving polynomials is to ask whether, given coefficients of a polynomial each of which is a symmetric function in the roots, one can "break" the symmetry and thereby recover the roots. Using this approach, solving a polynomial of degree is related to the ways of rearranging ("permuting") terms, called the symmetric group on letters and denoted . For the quadratic polynomial, the only ways to rearrange two roots are to either leave them be or to transpose them, so solving a quadratic polynomial is simple.
To find the roots and , consider their sum and difference:
These are called the Lagrange resolvents of the polynomial, from which the roots can be recovered as
Because is a symmetric function in and , it can be expressed in terms of and specifically as described above. However, is not symmetric, since exchanging and yields the additive inverse . So cannot be expressed in terms of the symmetric polynomials. However, its square is symmetric in the roots, expressible in terms of and . Specifically , which implies . Taking the positive root "breaks" the symmetry, resulting in
from which the roots and are recovered as
which is the quadratic formula for a monic polynomial.
Substituting , yields the usual expression for an arbitrary quadratic polynomial. The resolvents can be recognized as
respectively the vertex and the discriminant of the monic polynomial.
A similar but more complicated method works for cubic equations, which have three resolvents and a quadratic equation (the "resolving polynomial") relating and , which one can solve by the quadratic equation, and similarly for a quartic equation (degree 4), whose resolving polynomial is a cubic, which can in turn be solved. The same method for a quintic equation yields a polynomial of degree 24, which does not simplify the problem, and, in fact, solutions to quintic equations in general cannot be expressed using only roots.
Numerical calculation
The quadratic formula is exactly correct when performed using the idealized arithmetic of real numbers, but when approximate arithmetic is used instead, for example pen-and-paper arithmetic carried out to a fixed number of decimal places or the floating-point binary arithmetic available on computers, the limitations of the number representation can lead to substantially inaccurate results unless great care is taken in the implementation. Specific difficulties include catastrophic cancellation in computing the sum if ; catastrophic calculation in computing the discriminant itself in cases where ; degeneration of the formula when , , or is represented as zero or infinite; and possible overflow or underflow when multiplying or dividing extremely large or small numbers, even in cases where the roots can be accurately represented.
Catastrophic cancellation occurs when two numbers which are approximately equal are subtracted. While each of the numbers may independently be representable to a certain number of digits of precision, the identical leading digits of each number cancel, resulting in a difference of lower relative precision. When , evaluation of causes catastrophic cancellation, as does the evaluation of when . When using the standard quadratic formula, calculating one of the two roots always involves addition, which preserves the working precision of the intermediate calculations, while calculating the other root involves subtraction, which compromises it. Therefore, naïvely following the standard quadratic formula often yields one result with less relative precision than expected. Unfortunately, introductory algebra textbooks typically do not address this problem, even though it causes students to obtain inaccurate results in other school subjects such as introductory chemistry.
For example, if trying to solve the equation using a pocket calculator, the result of the quadratic formula might be approximately calculated as:
Even though the calculator used ten decimal digits of precision for each step, calculating the difference between two approximately equal numbers has yielded a result for with only four correct digits.
One way to recover an accurate result is to use the identity . In this example can be calculated as , which is correct to the full ten digits. Another more or less equivalent approach is to use the version of the quadratic formula with the square root in the denominator to calculate one of the roots (see above).
Practical computer implementations of the solution of quadratic equations commonly choose which formula to use for each root depending on the sign of .
These methods do not prevent possible overflow or underflow of the floating-point exponent in computing or , which can lead to numerically representable roots not being computed accurately. A more robust but computationally expensive strategy is to start with the substitution , turning the quadratic equation into
where is the sign function. Letting , this equation has the form , for which one solution is and the other solution is . The roots of the original equation are then and .
With additional complication the expense and extra rounding of the square roots can be avoided by approximating them as powers of two, while still avoiding exponent overflow for representable roots.
Historical development
The earliest methods for solving quadratic equations were geometric. Babylonian cuneiform tablets contain problems reducible to solving quadratic equations. The Egyptian Berlin Papyrus, dating back to the Middle Kingdom (2050 BC to 1650 BC), contains the solution to a two-term quadratic equation.
The Greek mathematician Euclid (circa 300 BC) used geometric methods to solve quadratic equations in Book 2 of his Elements, an influential mathematical treatise Rules for quadratic equations appear in the Chinese The Nine Chapters on the Mathematical Art circa 200 BC. In his work Arithmetica, the Greek mathematician Diophantus (circa 250 AD) solved quadratic equations with a method more recognizably algebraic than the geometric algebra of Euclid. His solution gives only one root, even when both roots are positive.
The Indian mathematician Brahmagupta included a generic method for finding one root of a quadratic equation in his treatise Brāhmasphuṭasiddhānta (circa 628 AD), written out in words in the style of the time but more or less equivalent to the modern symbolic formula. His solution of the quadratic equation was as follows: "To the absolute number multiplied by four times the [coefficient of the] square, add the square of the [coefficient of the] middle term; the square root of the same, less the [coefficient of the] middle term, being divided by twice the [coefficient of the] square is the value."
In modern notation, this can be written . The Indian mathematician Śrīdhara (8th–9th century) came up with a similar algorithm for solving quadratic equations in a now-lost work on algebra quoted by Bhāskara II. The modern quadratic formula is sometimes called Sridharacharya's formula in India and Bhaskara's formula in Brazil.
The 9th-century Persian mathematician Muḥammad ibn Mūsā al-Khwārizmī solved quadratic equations algebraically. The quadratic formula covering all cases was first obtained by Simon Stevin in 1594. In 1637 René Descartes published La Géométrie containing special cases of the quadratic formula in the form we know today.
Geometric significance
In terms of coordinate geometry, an axis-aligned parabola is a curve whose -coordinates are the graph of a second-degree polynomial, of the form , where , , and are real-valued constant coefficients with .
Geometrically, the quadratic formula defines the points on the graph, where the parabola crosses the -axis. Furthermore, it can be separated into two terms,
The first term describes the axis of symmetry, the line . The second term, , gives the distance the roots are away from the axis of symmetry.
If the parabola's vertex is on the -axis, then the corresponding equation has a single repeated root on the line of symmetry, and this distance term is zero; algebraically, the discriminant .
If the discriminant is positive, then the vertex is not on the -axis but the parabola opens in the direction of the -axis, crossing it twice, so the corresponding equation has two real roots. If the discriminant is negative, then the parabola opens in the opposite direction, never crossing the -axis, and the equation has no real roots; in this case the two complex-valued roots will be complex conjugates whose real part is the value of the axis of symmetry.
Dimensional analysis
If the constants , , and/or are not unitless then the quantities and must have the same units, because the terms and agree on their units. By the same logic, the coefficient must have the same units as , irrespective of the units of . This can be a powerful tool for verifying that a quadratic expression of physical quantities has been set up correctly.
| Mathematics | Elementary algebra | null |
59231 | https://en.wikipedia.org/wiki/Exception%20handling | Exception handling | In computing and computer programming, exception handling is the process of responding to the occurrence of exceptions – anomalous or exceptional conditions requiring special processing – during the execution of a program. In general, an exception breaks the normal flow of execution and executes a pre-registered exception handler; the details of how this is done depend on whether it is a hardware or software exception and how the software exception is implemented.
Exceptions are defined by different layers of a computer system, and the typical layers are CPU-defined interrupts, operating system (OS)-defined signals, programming language-defined exceptions. Each layer requires different ways of exception handling although they may be interrelated, e.g. a CPU interrupt could be turned into an OS signal. Some exceptions, especially hardware ones, may be handled so gracefully that execution can resume where it was interrupted.
Definition
The definition of an exception is based on the observation that each procedure has a precondition, a set of circumstances for which it will terminate "normally". An exception handling mechanism allows the procedure to raise an exception if this precondition is violated, for example if the procedure has been called on an abnormal set of arguments. The exception handling mechanism then handles the exception.
The precondition, and the definition of exception, is subjective. The set of "normal" circumstances is defined entirely by the programmer, e.g. the programmer may deem division by zero to be undefined, hence an exception, or devise some behavior such as returning zero or a special "ZERO DIVIDE" value (circumventing the need for exceptions). Common exceptions include an invalid argument (e.g. value is outside of the domain of a function), an unavailable resource (like a missing file, a network drive error, or out-of-memory errors), or that the routine has detected a normal condition that requires special handling, e.g., attention, end of file. Social pressure is a major influence on the scope of exceptions and use of exception-handling mechanisms, i.e. "examples of use, typically found in core libraries, and code examples in technical books, magazine articles, and online discussion forums, and in an organization’s code standards".
Exception handling solves the semipredicate problem, in that the mechanism distinguishes normal return values from erroneous ones. In languages without built-in exception handling such as C, routines would need to signal the error in some other way, such as the common return code and errno pattern. Taking a broad view, errors can be considered to be a proper subset of exceptions, and explicit error mechanisms such as errno can be considered (verbose) forms of exception handling. The term "exception" is preferred to "error" because it does not imply that anything is wrong - a condition viewed as an error by one procedure or programmer may not be viewed that way by another.
The term "exception" may be misleading because its connotation of "anomaly" indicates that raising an exception is abnormal or unusual, when in fact raising the exception may be a normal and usual situation in the program. For example, suppose a lookup function for an associative array throws an exception if the key has no value associated. Depending on context, this "key absent" exception may occur much more often than a successful lookup.
History
The first hardware exception handling was found in the UNIVAC I from 1951.
Arithmetic overflow executed two instructions at address 0 which could transfer control or fix up the result.
Software exception handling developed in the 1960s and 1970s. Exception handling was subsequently widely adopted by many programming languages from the 1980s onward.
Hardware exceptions
There is no clear consensus as to the exact meaning of an exception with respect to hardware. From the implementation point of view, it is handled identically to an interrupt: the processor halts execution of the current program, looks up the interrupt handler in the interrupt vector table for that exception or interrupt condition, saves state, and switches control.
IEEE 754 floating-point exceptions
Exception handling in the IEEE 754 floating-point standard refers in general to exceptional conditions and defines an exception as "an event that occurs when an operation on some particular operands has no outcome suitable for every reasonable application. That operation might signal one or more exceptions by invoking the default or, if explicitly requested, a language-defined alternate handling."
By default, an IEEE 754 exception is resumable and is handled by substituting a predefined value for different exceptions, e.g. infinity for a divide by zero exception, and providing status flags for later checking of whether the exception occurred (see C99 programming language for a typical example of handling of IEEE 754 exceptions). An exception-handling style enabled by the use of status flags involves: first computing an expression using a fast, direct implementation; checking whether it failed by testing status flags; and then, if necessary, calling a slower, more numerically robust, implementation.
The IEEE 754 standard uses the term "trapping" to refer to the calling of a user-supplied exception-handling routine on exceptional conditions, and is an optional feature of the standard. The standard recommends several usage scenarios for this, including the implementation of non-default pre-substitution of a value followed by resumption, to concisely handle removable singularities.
The default IEEE 754 exception handling behaviour of resumption following pre-substitution of a default value avoids the risks inherent in changing flow of program control on numerical exceptions. For example, the 1996 Cluster spacecraft launch ended in a catastrophic explosion due in part to the Ada exception handling policy of aborting computation on arithmetic error. William Kahan claims the default IEEE 754 exception handling behavior would have prevented this.
In programming languages
In user interfaces
Front-end web development frameworks, such as React and Vue, have introduced error handling mechanisms where errors propagate up the user interface (UI) component hierarchy, in a way that is analogous to how errors propagate up the call stack in executing code. Here the error boundary mechanism serves as an analogue to the typical try-catch mechanism. Thus a component can ensure that errors from its child components are caught and handled, and not propagated up to parent components.
For example, in Vue, a component would catch errors by implementing errorCapturedVue.component('parent', {
template: '<div><slot></slot></div>',
errorCaptured: (err, vm, info) => alert('An error occurred');
})
Vue.component('child', {
template: '<div>{{ cause_error() }}</div>'
})When used like this in markup:<parent>
<child></child>
</parent>The error produced by the child component is caught and handled by the parent component.
| Technology | Software development: General | null |
59238 | https://en.wikipedia.org/wiki/Entomology | Entomology | Entomology () is the branch of zoology that focuses on insects. Those who study entomology are known as entomologists. In the past, the term insect was less specific, and historically the definition of entomology would also include the study of animals in other arthropod groups, such as arachnids, myriapods, and crustaceans. The field is also referred to as insectology in American English, while in British English insectology implies the study of the relationships between insects and humans.
Over 1.3million insect species have been described by entomology.
History
Entomology is rooted in nearly all human cultures from prehistoric times, primarily in the context of agriculture (especially biological control and beekeeping). The natural Roman philosopher Pliny the Elder (23–79 CE) wrote a book on the kinds of insects, while the scientist of Kufa, Ibn al-A'rābī (760–845 CE) wrote a book on flies, (). However scientific study in the modern sense began only relatively recently, in the 16th century. Ulisse Aldrovandi's (Concerning Insect Animals) was published in 1602. Microscopist Jan Swammerdam published History of Insects, correctly describing the reproductive organs of insects and metamorphosis. In 1705, Maria Sibylla Merian published the book about the tropical insects of Dutch Surinam.
Early entomological works associated with the naming and classification of species followed the practice of maintaining cabinets of curiosity, predominantly in Europe. This collecting fashion led to the formation of natural history societies, exhibitions of private collections, and journals for recording communications and the documentation of new species. Many of the collectors tended to be from the aristocracy, and there developed a trade involving collectors around the world and traders. This has been called the "era of heroic entomology". William Kirby is widely considered as the father of entomology in England. In collaboration with William Spence, he published a definitive entomological encyclopedia, Introduction to Entomology, regarded as the subject's foundational text. He also helped found the Royal Entomological Society in London in 1833, one of the earliest such societies in the world; earlier antecedents, such as the Aurelian society date back to the 1740s. In the late 19th century, the growth of agriculture, and colonial trade spawned the "era of economic entomology" which created the professional entomologist associated with the rise of the university and training in the field of biology.
Entomology developed rapidly in the 19th and 20th centuries and was studied by large numbers of people, including such notable figures as Charles Darwin, Jean-Henri Fabre, Vladimir Nabokov, Karl von Frisch (winner of the 1973 Nobel Prize in Physiology or Medicine), and twice Pulitzer Prize winner E. O. Wilson.
There has also been a history of people becoming entomologists through museum curation and research assistance, such as Sophie Lutterlough at the Smithsonian National Museum of Natural History. Insect identification is an increasingly common hobby, with butterflies and (to a lesser extent) dragonflies being the most popular.
Most insects can easily be allocated to order, such as Hymenoptera (bees, wasps, and ants) or Coleoptera (beetles). However, identifying to genus or species is usually only possible through the use of identification keys and monographs. Because the class Insecta contains a very large number of species (over 330,000 species of beetles alone) and the characteristics distinguishing them are unfamiliar, and often subtle (or invisible without a microscope), this is often very difficult even for a specialist. This has led to the development of automated species identification systems targeted on insects, for example, Daisy, ABIS, SPIDA and Draw-wing.
Applications
Pest control
In 1994, the Entomological Society of America launched a new professional certification program for the pest control industry called the Associate Certified Entomologist (ACE). To qualify as a "true entomologist" an individual would normally require an advanced degree, with most entomologists pursuing a PhD. While not true entomologists in the traditional sense, individuals who attain the ACE certification may be referred to as ACEs or Associate Certified Entomologists.
As such, other credential programs managed by the Entomological Society of America have varying credential requirements. These different programs are known as Public Health Entomology (PHE), Certified IPM Technicians (CITs), and Board Certified Entomologists (BCEs) (ESA Certification Corporation). To be qualified in public health entomology (PHE), one must pass an exam on the types of arthropods that can spread diseases and lead to medical complications (ESA Certification Corporation). These individuals also have to "agree to ascribe to a code of ethical behavior" (ESA Certification Corporation). Individuals who are planning to become Certified IPM Technicians (CITs), need to obtain at around 1-4 years of experience in pest management and successfully pass an exam, that is based on the information, that they are acquainted with (ESA Certification Corporation). Like in Public Health Entomology (PHE), those who want to become Certified IPM Technicians (CITs) also have to "agree to ascribe to a code of ethical behavior" (ESA Certification Corporation). These individuals must also be approved to use pesticides (ESA Certification Corporation). For those who plan on becoming Board Certified Entomologists (BCEs), individuals have to pass two exams and "agree to ascribe to a code of ethical behavior" (ESA Certification Corporation). As with this, they also have to fulfill a certain amount of educational requirements every 12 months (ESA Certification Corporation).
Forensics
Forensic entomology is a branch of forensic science that studies insects found on corpses or elsewhere around crime scenes. This includes studying the types of insects commonly found on cadavers, their life cycles, their presence in different environments, and how insect assemblages change with decomposition.
Medicine
Medical entomology is focused upon insects and arthropods that impact human health. Veterinary entomology is included in this category, because many animal diseases can "jump species" and become a human health threat, for example, bovine encephalitis. Medical entomology also includes scientific research on the behavior, ecology, and epidemiology of arthropod disease vectors, and involves a tremendous outreach to the public, including local and state officials and other stake holders in the interest of public safety.
Subdisciplines
Many entomologists specialize in a single order or even a family of insects, and a number of these subspecialties are given their own generic names, typically (but not always) derived from the scientific name of the group:
Coleopterology – beetles
Dipterology – flies
Odonatology – dragonflies and damselflies
Hemipterology – true bugs
Isopterology – termites
Lepidopterology – moths and butterflies
Melittology (or Apiology) – bees
Myrmecology – ants
Orthopterology – grasshoppers, crickets, etc.
Trichopterology – caddisflies
Vespology – social wasps
Organizations
Like other scientific specialties, entomologists have a number of local, national, and international organizations. There are also many organizations specializing in specific subareas.
Amateur Entomologists' Society
British Entomological and Natural History Society
Entomological Society of America
Entomological Society of Canada
Entomological Society of Japan
Entomologischer Verein Krefeld
Royal Entomological Society
Australian Entomological Society
Entomological Society of New Zealand
Research collection
Here is a list of selected very large insect collections, housed in museums, universities, or research institutes.
Asia
Zoological Survey of India
Africa
Natal Museum, Pietermaritzburg, South Africa
Australasia
Lincoln University Entomology Research Collection, Lincoln, New Zealand
Europe
Bavarian State Collection of Zoology, Zoologische Staatssammlung München
United States
Academy of Natural Sciences of Philadelphia
Canada
Canadian Museum of Nature, Ottawa, Ontario
| Biology and health sciences | Basics_2 | Biology |
59243 | https://en.wikipedia.org/wiki/Amphibole | Amphibole | Amphibole ( ) is a group of inosilicate minerals, forming prism or needlelike crystals, composed of double chain tetrahedra, linked at the vertices and generally containing ions of iron and/or magnesium in their structures. Its IMA symbol is Amp. Amphiboles can be green, black, colorless, white, yellow, blue, or brown. The International Mineralogical Association currently classifies amphiboles as a mineral supergroup, within which are two groups and several subgroups.
Mineralogy
Amphiboles crystallize into two crystal systems, monoclinic and orthorhombic. In chemical composition and general characteristics they are similar to the pyroxenes. The chief differences from pyroxenes are that (i) amphiboles contain essential hydroxyl (OH) or halogen (F, Cl) and (ii) the basic structure is a double chain of tetrahedra (as opposed to the single chain structure of pyroxene). Most apparent, in hand specimens, is that amphiboles form oblique cleavage planes (at around 120 degrees), whereas pyroxenes have cleavage angles of approximately 90 degrees. Amphiboles are also specifically less dense than the corresponding pyroxenes. Amphiboles are the primary constituent of amphibolites.
Structure
Like pyroxenes, amphiboles are classified as inosilicate (chain silicate) minerals. However, the pyroxene structure is built around single chains of silica tetrahedra while amphiboles are built around double chains of silica tetrahedra. In other words, as with almost all silicate minerals, each silicon ion is surrounded by four oxygen ions. In amphiboles, some of the oxygen ions are shared between silicon ions to form a double chain structure as depicted below. These chains extend along the [001] axis of the crystal. One side of each chain has apical oxygen ions, shared by only one silicon ion, and pairs of double chains are bound to each other by metal ions that connect apical oxygen ions. The pairs of double chains have been likened to I-beams. Each I-beam is bonded to its neighbor by additional metal ions to form the complete crystal structure. Large gaps in the structure may be empty or partially filled by large metal ions, such as sodium, but remain points of weakness that help define the cleavage planes of the crystal.
In rocks
Amphiboles are minerals of either igneous or metamorphic origin. Amphiboles are more common in intermediate to felsic igneous rocks than in mafic igneous rocks, because the higher silica and dissolved water content of the more evolved magmas favors formation of amphiboles rather than pyroxenes. The highest amphibole content, around 20%, is found in andesites. Hornblende is widespread in igneous and metamorphic rocks and is particularly common in syenites and diorites. Calcium is sometimes a constituent of naturally occurring amphiboles. Amphiboles of metamorphic origin include those developed in limestones by contact metamorphism (tremolite) and those formed by the alteration of other ferromagnesian minerals (such as hornblende as an alteration product of pyroxene). Pseudomorphs of amphibole after pyroxene are known as uralite.
History and etymology
The name amphibole derives from Greek (, ), implying ambiguity. The name was used by to include tremolite, actinolite and hornblende. The group was so named by Haüy in allusion to the protean variety, in composition and appearance, assumed by its minerals. This term has since been applied to the whole group. Numerous sub-species and varieties are distinguished, the more important of which are tabulated below in two series. The formulae of each will be seen to be built on the general double-chain silicate formula RSi4O11.
Four of the amphibole minerals are commonly called asbestos. These are: anthophyllite, riebeckite, the cummingtonite/grunerite series, and the actinolite/tremolite series. The cummingtonite/grunerite series is often termed amosite or "brown asbestos", and riebeckite is known as crocidolite or "blue asbestos". These are generally called amphibole asbestos. Mining, manufacture and prolonged use of these minerals can cause serious illnesses.
Mineral species
The more common amphiboles are classified as shown in the following table:
Other species
Orthorhombic series
Holmquistite, Li2Mg3Al2Si8O22(OH)2
Monoclinic series
Pargasite, NaCa2Mg3Fe2+Si6Al3O22(OH)2
Winchite, (CaNa)Mg4(Al,Fe3+)Si8O22(OH)2
Edenite, NaCa2Mg5(Si7Al)O22(OH)2
Series
Certain amphibole minerals form solid solution series, at least at elevated temperature. Ferrous iron usually substitutes freely for magnesium in amphiboles to form continuous solid solution series between magnesium-rich and iron-rich endmembers. These include the cummington (magnesium) to grunerite (iron) endmembers, where the dividing line is placed at 30% magnesium.
In addition, the orthoamphiboles, anthophyllite and gedrite, which differ in their aluminium content, form a continuous solid solution at elevated temperature. As the amphibole cools, the two end members exsolve to form very thin layers (lamellae).
Hornblende is highly variable in composition, and includes at least five solid solution series: magnesiohornblende-ferrohornblende (), tschermakite-ferrotschermakite (), edenite-ferroedenite (), pargasite-ferropargasite () and magnesiohastingstite-hastingsite (). In addition, titanium, manganese, or chromium can substitute for some of the cations and oxygen, fluorine, or chlorine for some of the hydroxide. The different chemical types are almost impossible to distinguish even by optical or X-ray methods, and detailed chemical analysis using an electron microprobe is required.
Glaucophane to riebeckite form yet another solid solution series, which also extends towards hornblende and arfvedsonite.
There is not a continuous series between calcic clinoamphiboles, such as hornblende, and low-calcium amphiboles, such as orthoamphiboles or the cummingtonite-grunerite series. Compositions intermediate in calcium are almost nonexistent in nature. However, there is a solid solution series between hornblende and tremolite-actinolite at elevated temperature. A miscibility gap exists at lower temperatures, and, as a result, hornblende often contains exsolution lamellae of grunerite.
Descriptions
On account of the wide variations in chemical composition, the different members vary considerably in properties and general appearance.
Anthophyllite occurs as brownish, fibrous or lamellar masses with hornblende in mica-schist at Kongsberg in Norway and some other localities. An aluminous related species is known as gedrite and a deep green Russian variety containing little iron as kupfferite.
Hornblende is an important constituent of many igneous rocks. It is also an important constituent of amphibolites formed by metamorphism of basalt.
Actinolite is an important and common member of the monoclinic series, forming radiating groups of acicular crystals of a bright green or greyish-green color. It occurs frequently as a constituent of greenschists. The name (from Greek ἀκτίς, ἀκτῖνος/aktís, aktînos, a 'ray' and λίθος/líthos, a 'stone') is a translation of the old German word Strahlstein (radiated stone).
Glaucophane, crocidolite, riebeckite and arfvedsonite form a somewhat special group of alkali-amphiboles. The first two are blue fibrous minerals, with glaucophane occurring in blueschists and crocidolite (blue asbestos) in ironstone formations, both resulting from dynamo-metamorphic processes. The latter two are dark green minerals, which occur as original constituents of igneous rocks rich in sodium, such as nepheline-syenite and phonolite.
Pargasite is a rare magnesium-rich variety of hornblende with essential sodium, usually found in ultramafic rocks. For instance, it occurs in uncommon mantle xenoliths, carried up by kimberlite. It is hard, dense, black and usually automorphic, with a red-brown pleochroism in petrographic thin section.
| Physical sciences | Silicate minerals | Earth science |
59249 | https://en.wikipedia.org/wiki/Papaya | Papaya | The papaya (, ), papaw, () or pawpaw () is the plant species Carica papaya, one of the 21 accepted species in the genus Carica of the family Caricaceae, and also the name of its fruit. It was first domesticated in Mesoamerica, within modern-day southern Mexico and Central America. It is grown in several countries in regions with a tropical climate. In 2022, India produced 38% of the world's supply of papayas.
Etymology
The word papaya derives from Arawak via Spanish, and is also the name for the plant. The name papaw or pawpaw is used alternatively for the fruit only in some regions.
Description
The papaya is a small, sparsely branched tree, usually with a single stem growing from tall, with spirally arranged leaves confined to the top of the trunk. The lower trunk is conspicuously scarred where leaves and fruit were borne. The leaves are large, in diameter, deeply palmately lobed, with seven lobes. All plant parts contain latex in articulated laticifers.
Flowers
Papayas are dioecious. The flowers are five-parted and highly dimorphic; the male flowers have the stamens fused to the petals. There are two different types of papaya flowers. The female flowers have a superior ovary and five contorted petals loosely connected at the base.
Male and female flowers are borne in the leaf axils; the male flowers are in multiflowered dichasia, and the female ones are in few-flowered dichasia. The pollen grains are elongated and approximately 35 microns in length. The flowers are sweet-scented, open at night, and are wind- or insect-pollinated.
Fruit
The fruit is a large berry about long and in diameter. It is ripe when it feels soft (as soft as a ripe avocado or softer) and its skin has attained an amber to orange hue. Along the walls of the large central cavity are attached numerous black seeds.
Chemistry
Papaya skin, pulp, and seeds contain a variety of phytochemicals, including carotenoids and polyphenols, as well as benzyl isothiocyanates and benzyl glucosinates, with skin and pulp levels that increase during ripening. The carotenoids, lutein and beta-carotene, are prominent in the yellow skin, while lycopene is dominant in the red flesh (table). Papaya seeds also contain the cyanogenic substance prunasin. The green fruit contains papain, a cysteine protease enzyme used to tenderize meat (see below).
Distribution and habitat
Native to tropical America, papaya originates from southern Mexico and Central America. Papaya is also considered native to southern Florida, introduced by predecessors of the Calusa no later than AD 300. Spaniards introduced papaya to the Old World in the 16th century. Papaya cultivation is now nearly pantropical, spanning Hawaii, Central Africa, India, and Australia.
Wild populations of papaya are generally confined to naturally disturbed tropical forests. Papaya is found in abundance on Everglades hammocks following major hurricanes, but is otherwise infrequent. In the rain forests of southern Mexico, papaya thrives and reproduces quickly in canopy gaps while dying off in the mature closed-canopy forests.
Ecology
Viruses
Papaya ringspot virus is a well-known virus within plants in Florida. The first signs of the virus are yellowing and vein-clearing of younger leaves and mottling yellow leaves. Infected leaves may obtain blisters, roughen, or narrow, with blades sticking upwards from the middle of the leaves. The petioles and stems may develop dark green greasy streaks and, in time, become shorter. The ringspots are circular, C-shaped markings that are a darker green than the fruit. In the later stages of the virus, the markings may become gray and crusty. Viral infections impact growth and reduce the fruit's quality. One of the biggest effects that viral infections have on papaya is taste. As of 2010, the only way to protect papaya from this virus is genetic modification.
The papaya mosaic virus destroys the plant until only a small tuft of leaves is left. The virus affects both the leaves of the plant and the fruit. Leaves show thin, irregular, dark-green lines around the borders and clear areas around the veins. The more severely affected leaves are irregular and linear in shape. The virus can infect the fruit at any stage of its maturity. Fruits as young as two weeks old have been spotted with dark-green ringspots about 1 inch (25 mm) in diameter. Rings on the fruit are most likely seen on either the stem end or the blossom end. In the early stages of the ringspots, the rings tend to be many closed circles, but as the disease develops, the rings increase in diameter consisting of one large ring. The difference between the ringspot and the mosaic viruses is the ripe fruit in the ringspot has a mottling of colors, and the mosaic does not.
Fungi and oomycetes
The fungus anthracnose attacks papaya, especially mature fruits. The disease starts small with very few signs, such as water-soaked spots on ripening fruits. The spots become sunken, turn brown or black, and may get bigger. In some of the older spots, the fungus may produce pink spores. The fruit ends up being soft and having an off flavor because the fungus grows into the fruit.
The fungus powdery mildew occurs as a superficial white presence on the leaf's surface, which is easily recognized. Tiny, light yellow spots begin on the lower surfaces of the leaf as the disease starts to make its way. The spots enlarge, and white powdery growth appears on the leaves. The infection usually appears at the upper leaf surface as white fungal growth. Powdery mildew is not as severe as other diseases.
The fungus-like oomycete Phytophthora causes damping-off, root rot, stem rot, stem girdling, and fruit rot. Damping-off happens in young plants by wilting and death. The spots on established plants start as white, water-soaked lesions at the fruit and branch scars. These spots enlarge and eventually cause death. The disease's most dangerous feature is the fruit's infection, which may be toxic to consumers. The roots can also be severely and rapidly infected, causing the plant to brown and wilt away, collapsing within days.
Pests
The papaya fruit fly lays its eggs inside of the fruit, possibly up to 100 or more eggs. The eggs usually hatch within 12 days when they begin to feed on seeds and interior parts of the fruit. When the larvae mature, usually 16 days after being hatched, they eat their way out of the fruit, drop to the ground, and pupate in the soil to emerge within one to two weeks later as mature flies. The infected papaya turns yellow and drops to the ground after the papaya fruit fly infestation.
The two-spotted spider mite is a 0.5-mm-long brown or orange-red or a green, greenish-yellow translucent oval pest. They all have needle-like piercing-sucking mouthparts and feed by piercing the plant tissue with their mouthparts, usually on the underside of the plant. The spider mites spin fine threads of webbing on the host plant, and when they remove the sap, the mesophyll tissue collapses, and a small chlorotic spot forms at the feeding sites. The leaves of the papaya fruit turn yellow, gray, or bronze. If the spider mites are not controlled, they can cause the death of the fruit.
The papaya whitefly lays yellow, oval eggs that appear dusted on the undersides of the leaves. They eat papaya leaves, therefore damaging the fruit. There, the eggs developed into flies in three stages called instars. The first instar has well-developed legs and is the only mobile immature life stage. The crawlers insert their mouthparts in the lower surfaces of the leaf when they find it suitable and usually do not move again in this stage. The next instars are flattened, oval, and scale-like. In the final stage, the pupal whiteflies are more convex, with large, conspicuously red eyes.
Papayas are one of the most common hosts for fruit flies like A. suspensa, which lay their eggs in overripe or spoiled papayas. The larvae of these flies then consume the fruit to gain nutrients until they can proceed into the pupal stage. This parasitism has led to extensive economic costs for nations in Central America.
Cultivation
Historical accounts from 18th-century travelers and botanists suggested that papaya seeds were transported from the Caribbean to Malacca and then to India. From Malacca or the Philippines, papaya spread throughout Asia and into the South Pacific region. Credit for introducing papaya to Hawaii is often given to Francisco de Paula Marín, a Spanish explorer and horticulturist, who brought it from the Marquesas Islands in the early 1800s. Since then, papaya cultivation has expanded to all tropical countries and many subtropical regions worldwide. Today, papaya is grown extensively across the globe, owing to its adaptability to various climates and its popularity as a tropical fruit.
Papaya plants grow in three sexes: male, female, and hermaphrodite. The male produces only pollen, never fruit. The female produces small, inedible fruits unless pollinated. The hermaphrodite can self-pollinate since its flowers contain both male stamens and female ovaries. Almost all commercial papaya orchards contain only hermaphrodites.
Originally from southern Mexico (particularly Chiapas and Veracruz), Central America, northern South America, and southern Florida the papaya is now cultivated in most tropical countries. In cultivation, it grows rapidly, fruiting within three years. It is, however, highly frost-sensitive, limiting its production to tropical climates. Temperatures below are greatly harmful, if not fatal. In Florida, California, and Texas, growth is generally limited to the southern parts of those states. It prefers sandy, well-drained soil, as standing water can kill the plant within 24 hours.
Cultivars
Two kinds of papayas are commonly grown. One has sweet, red, or orange flesh, and the other has yellow flesh; in Australia, these are called "red papaya" and "yellow papaw," respectively. Either kind, picked green, is called a "green papaya."
The large-fruited, red-fleshed 'Maradol,' 'Sunrise,' and 'Caribbean Red' papayas often sold in U.S. markets are commonly grown in Mexico and Belize.
In 2011, Philippine researchers reported that by hybridizing papaya with Vasconcellea quercifolia, they had developed papaya resistant to papaya ringspot virus (PRV), part of a long line of attempts to transfer resistance from Vasconcellea species into papaya.
Genetically engineered cultivars
Carica papaya was the first transgenic fruit tree to have its genome sequenced. In response to the papaya ringspot virus outbreak in Hawaii in 1998, genetically altered papaya were approved and brought to market (including 'SunUp' and 'Rainbow' varieties.) Varieties resistant to PRV have some DNA of this virus incorporated into the plant's DNA. As of 2010, 80% of Hawaiian papaya plants were genetically modified. The modifications were made by University of Hawaii scientists, who made the modified seeds available to farmers without charge.
In transgenic papaya, resistance is produced by inserting the viral coat protein gene into the plant's genome. Doing so seems to cause a similar protective reaction in the plant to cross-protection, which involves using an attenuated virus to protect against a more dangerous strain. Conventional varieties of transgenic papaya has reduced resistance against heterologous (not closely related to the coat gene source) strains, forcing different localities to develop their own transgenic varieties. As of 2016, one transgenic line appears able to deal with three different heterologous strains in addition to its source.
Production
In 2022, global production of papayas was 13.8 million tonnes, led by India with 38% of the world total (table). Global papaya production grew significantly over the early 21st century, mainly as a result of increased production in India and demand by the United States. The United States is the largest importer of papayas worldwide.
In South Africa, papaya orchards yield up to 100 tonnes of fruit per hectare.
Toxicity
Papaya releases a latex fluid when not ripe, possibly causing irritation and an allergic reaction in some people. Because the enzyme papain acts as an allergen in sensitive individuals, meat that has been tenderized with it may induce an allergic reaction.
Culinary use
The ripe fruit of the papaya is usually eaten raw, without skin or seeds. The black seeds are edible and have a sharp, spicy taste. The unripe green fruit is usually cooked due to its latex content.
Both green papaya fruit and its latex are rich in papain, a cysteine protease used for tenderizing meat and other proteins, as practiced currently by indigenous Americans, people of the Caribbean region, Pacific Islands, and the Philippines. It is included as a component in some powdered meat tenderizers. Papaya is not suitable for foods which set due to gelatin (such as jelly or aspic) because the enzymatic properties of papain prevent gelatin from setting.
Nutrition
Raw papaya pulp is 88% water, 11% carbohydrates, and contains negligible fat and protein (table). In a reference amount of , papaya fruit provides 43 kilocalories and is a significant source of vitamin C (69% of the Daily Value, DV) and a moderate source of folate (10% DV), but otherwise has a low content of micronutrients (table).
Southeast Asia
Green papaya is used in Southeast Asian cooking, both raw and cooked. In some parts of Asia, the young leaves of the papaya are steamed and eaten like spinach.
In Myanmar, the unripe papaya are cut into slices and dipped into sour, fermented, or spicy seasonings and dips. In Myanmar and Thai recipes, the unripe papaya are cut into thinner slices to make papaya salad. The reason the unripe papaya is used is because of the firmer and crunchier texture.
Papayas became a part of Filipino cuisine after being introduced to the islands via the Manila galleons. Unripe or nearly ripe papayas (with orange flesh but still hard and green) are julienned and are commonly pickled into atchara, which is ubiquitous as a side dish to salty dishes. Nearly ripe papayas can also be eaten fresh as ensaladang papaya (papaya salad) or cubed and eaten dipped in vinegar or salt. Green papaya is also a common ingredient or filling in various savory dishes such as okoy, tinola, ginataan, lumpia, and empanada, especially in the cuisines of northern Luzon.
In Indonesian cuisine, the unripe green fruits and young leaves are boiled for use as part of lalab salad, while the flower buds are sautéed and stir-fried with chilies and green tomatoes as Minahasan papaya flower vegetable dish.
In Lao and Thai cuisine, unripe green papayas are used to make a type of spicy salad known in Laos as tam maak hoong and in Thailand as som tam. It is also used in Thai curries, such as kaeng som.
South America
In Brazil and Paraguay, the unripe fruits are used to make sweets or preserves.
Traditional medicine
In traditional medicine, papaya leaves have been believed useful as a treatment for malaria, an abortifacient, a purgative, or smoked to relieve asthma.
| Biology and health sciences | Brassicales | null |
59358 | https://en.wikipedia.org/wiki/Mycorrhiza | Mycorrhiza | A mycorrhiza (; , mycorrhiza, or mycorrhizas) is a symbiotic association between a fungus and a plant. The term mycorrhiza refers to the role of the fungus in the plant's rhizosphere, the plant root system and its surroundings. Mycorrhizae play important roles in plant nutrition, soil biology, and soil chemistry.
In a mycorrhizal association, the fungus colonizes the host plant's root tissues, either intracellularly as in arbuscular mycorrhizal fungi, or extracellularly as in ectomycorrhizal fungi. The association is normally mutualistic. In particular species, or in particular circumstances, mycorrhizae may have a parasitic association with host plants.
Definition
A mycorrhiza is a symbiotic association between a green plant and a fungus. The plant makes organic molecules by photosynthesis and supplies them to the fungus in the form of sugars or lipids, while the fungus supplies the plant with water and mineral nutrients, such as phosphorus, taken from the soil. Mycorrhizas are located in the roots of vascular plants, but mycorrhiza-like associations also occur in bryophytes and there is fossil evidence that early land plants that lacked roots formed arbuscular mycorrhizal associations. Most plant species form mycorrhizal associations, though some families like Brassicaceae and Chenopodiaceae cannot. Different forms for the association are detailed in the next section. The most common is the arbuscular type that is present in 70% of plant species, including many crop plants such as cereals and legumes.
Evolution
Fossil and genetic evidence indicate that mycorrhizae are ancient, potentially as old as the terrestrialization of plants. Genetic evidence indicates that all land plants share a single common ancestor, which appears to have quickly adopted mycorrhizal symbiosis, and research suggests that proto-mycorrhizal fungi were a key factor enabling plant terrestrialization. The 400 million year old Rhynie chert contains an assemblage of fossil plants preserved in sufficient detail that arbuscular mycorrhizae have been observed in the stems of Aglaophyton major, giving a lower bound for how late mycorrhizal symbiosis may have developed. Ectomycorrhizae developed substantially later, during the Jurassic period, while most other modern mycorrhizal families, including orchid and ericoid mycorrhizae, date to the period of angiosperm radiation in the Cretaceous period. There is genetic evidence that the symbiosis between legumes and nitrogen-fixing bacteria is an extension of mycorrhizal symbiosis. The modern distribution of mycorrhizal fungi appears to reflect an increasing complexity and competition in root morphology associated with the dominance of angiosperms in the Cenozoic Era, characterized by complex ecological dynamics between species.
Types
The mycorrhizal lifestyle has independently convergently evolved multiple times in the history of Earth. There are multiple ways to categorize mycorrhizal symbiosis. One major categorization is the division between ectomycorrhizas and endomycorrhizas. The two types are differentiated by the fact that the hyphae of ectomycorrhizal fungi do not penetrate individual cells within the root, while the hyphae of endomycorrhizal fungi penetrate the cell wall and invaginate the cell membrane.
Similar symbiotic relationships
Some forms of plant-fungal symbiosis are similar to mycorrhizae, but considered distinct. One example is fungal endophytes. Endophytes are defined as organisms that can live within plant cells without causing harm to the plant. They are distinguishable from mycorrhizal fungi by the absence of nutrient-transferring structures for bringing in nutrients from outside the plant. Some lineages of mycorrhizal fungi may have evolved from endophytes into mycorrhizal fungi, and some fungi can live as mycorrhizae or as endophytes.
Ectomycorrhiza
Ectomycorrhizae are distinct in that they do not penetrate into plant cells, but instead form a structure called a Hartig net that penetrates between cells. Ectomycorrhizas consist of a hyphal sheath, or mantle, covering the root tip and the Hartig net of hyphae surrounding the plant cells within the root cortex. In some cases the hyphae may also penetrate the plant cells, in which case the mycorrhiza is called an endomycorrhiza. Outside the root, ectomycorrhizal extramatrical mycelium forms an extensive network within the soil and leaf litter. Other forms of mycorrhizae, including arbuscular, ericoid, arbutoid, monotropoid, and orchid mycorrhizas, are considered endomycorrhizae.
Ectomycorrhizas, or EcM, are symbiotic associations between the roots of around 10% of plant families, mostly woody plants including the birch, dipterocarp, eucalyptus, oak, pine, and rose families, orchids, and fungi belonging to the Basidiomycota, Ascomycota, and Zygomycota. Ectomycorrhizae associate with relatively few plant species, only about 2% of plant species on Earth, but the species they associate with are mostly trees and woody plants that are highly dominant in their ecosystems, meaning plants in ectomycorrhizal relationships make up a large proportion of plant biomass. Some EcM fungi, such as many Leccinum and Suillus, are symbiotic with only one particular genus of plant, while other fungi, such as the Amanita, are generalists that form mycorrhizas with many different plants. An individual tree may have 15 or more different fungal EcM partners at one time. While the diversity of plants involved in EcM is low, the diversity of fungi involved in EcM is high. Thousands of ectomycorrhizal fungal species exist, hosted in over 200 genera. A recent study has conservatively estimated global ectomycorrhizal fungal species richness at approximately 7750 species, although, on the basis of estimates of knowns and unknowns in macromycete diversity, a final estimate of ECM species richness would probably be between 20,000 and 25,000. Ectomycorrhizal fungi evolved independently from saprotrophic ancestors many times in the group's history.
Nutrients can be shown to move between different plants through the fungal network. Carbon has been shown to move from paper birch seedlings into adjacent Douglas-fir seedlings, although not conclusively through a common mycorrhizal network, thereby promoting succession in ecosystems. The ectomycorrhizal fungus Laccaria bicolor has been found to lure and kill springtails to obtain nitrogen, some of which may then be transferred to the mycorrhizal host plant. In a study by Klironomos and Hart, Eastern White Pine inoculated with L. bicolor was able to derive up to 25% of its nitrogen from springtails. When compared with non-mycorrhizal fine roots, ectomycorrhizae may contain very high concentrations of trace elements, including toxic metals (cadmium, silver) or chlorine.
The first genomic sequence for a representative of symbiotic fungi, the ectomycorrhizal basidiomycete L. bicolor, was published in 2008. An expansion of several multigene families occurred in this fungus, suggesting that adaptation to symbiosis proceeded by gene duplication. Within lineage-specific genes those coding for symbiosis-regulated secreted proteins showed an up-regulated expression in ectomycorrhizal root tips suggesting a role in the partner communication. L. bicolor is lacking enzymes involved in the degradation of plant cell wall components (cellulose, hemicellulose, pectins and pectates), preventing the symbiont from degrading host cells during the root colonisation. By contrast, L. bicolor possesses expanded multigene families associated with hydrolysis of bacterial and microfauna polysaccharides and proteins. This genome analysis revealed the dual saprotrophic and biotrophic lifestyle of the mycorrhizal fungus that enables it to grow within both soil and living plant roots. Since then, the genomes of many other ectomycorrhizal fungal species have been sequenced further expanding the study of gene families and evolution in these organisms.
Arbutoid mycorrhiza
This type of mycorrhiza involves plants of the Ericaceae subfamily Arbutoideae. It is however different from ericoid mycorrhiza and resembles ectomycorrhiza, both functionally and in terms of the fungi involved. It differs from ectomycorrhiza in that some hyphae actually penetrate into the root cells, making this type of mycorrhiza an ectendomycorrhiza.
Arbuscular mycorrhiza
Arbuscular mycorrhizas, (formerly known as vesicular-arbuscular mycorrhizas), have hyphae that penetrate plant cells, producing branching, tree-like structures called arbuscules within the plant cells for nutrient exchange. Often, balloon-like storage structures, termed vesicles, are also produced. In this interaction, fungal hyphae do not in fact penetrate the protoplast (i.e. the interior of the cell), but invaginate the cell membrane, creating a so-called peri-arbuscular membrane. The structure of the arbuscules greatly increases the contact surface area between the hypha and the host cell cytoplasm to facilitate the transfer of nutrients between them. Arbuscular mycorrhizas are obligate biotrophs, meaning that they depend upon the plant host for both growth and reproduction; they have lost the ability to sustain themselves by decomposing dead plant material. Twenty percent of the photosynthetic products made by the plant host are consumed by the fungi, the transfer of carbon from the terrestrial host plant is then exchanged by equal amounts of phosphate from the fungi to the plant host.
Contrasting with the pattern seen in ectomycorrhizae, the species diversity of AMFs is very low, but the diversity of plant hosts is very high; an estimated 78% of all plant species associate with AMFs. Arbuscular mycorrhizas are formed only by fungi in the division Glomeromycota. Fossil evidence and DNA sequence analysis suggest that this mutualism appeared 400-460 million years ago, when the first plants were colonizing land. Arbuscular mycorrhizas are found in 85% of all plant families, and occur in many crop species. The hyphae of arbuscular mycorrhizal fungi produce the glycoprotein glomalin, which may be one of the major stores of carbon in the soil. Arbuscular mycorrhizal fungi have (possibly) been asexual for many millions of years and, unusually, individuals can contain many genetically different nuclei (a phenomenon called heterokaryosis).
Mucoromycotina fine root endophytes
Mycorrhizal fungi belonging to Mucoromycotina, known as “fine root endophytes" (MFREs), were mistakenly identified as arbuscular mycorrhizal fungi until recently. While similar to AMF, MFREs are from subphylum Mucoromycotina instead of Glomeromycotina. Their morphology when colonizing a plant root is very similar to AMF, but they form fine textured hyphae. Effects of MFREs may have been mistakenly attributed to AMFs due to confusion between the two, complicated by the fact that AMFs and MFREs often colonize the same hosts simultaneously. Unlike AMFs, they appear capable of surviving without a host. This group of mycorrhizal fungi is little understood, but appears to prefer wet, acidic soils and forms symbiotic relationships with liverworts, hornworts, lycophytes, and angiosperms.
Ericoid mycorrhiza
Ericoid mycorrhizae, or ErMs, involve only plants in Ericales and are the most recently evolved of the major mycorrhizal relationships. Plants that form ericoid mycorrhizae are mostly woody understory shrubs; hosts include blueberries, bilberries, cranberries, mountain laurels, rhododendrons, heather, neinei, and giant grass tree. ErMs are most common in boreal forests, but are found in two-thirds of all forests on Earth. Ericoid mycorrhizal fungi belong to several different lineages of fungi. Some species can live as endophytes entirely within plant cells even within plants outside the Ericales, or live independently as saprotrophs that decompose dead organic matter. This ability to switch between multiple lifestyle types makes ericoid mycorrhizal fungi very adaptable.
Plants that participate in these symbioses have specialized roots with no root hairs, which are covered with a layer of epidermal cells that the fungus penetrates into and completely occupies. The fungi have a simple intraradical (growth in cells) phase, consisting of dense coils of hyphae in the outermost layer of root cells. There is no periradical phase and the extraradical phase consists of sparse hyphae that don't extend very far into the surrounding soil. They might form sporocarps (probably in the form of small cups), but their reproductive biology is poorly understood.
Plants participating in ericoid mycorrhizal symbioses are found in acidic, nutrient-poor conditions. Whereas AMFs have lost their saprotrophic capabilities, and EcM fungi have significant variation in their ability to produce enzymes needed for a saprotrophic lifestyle, fungi involved in ErMs have fully retained the ability to decompose plant material for sustenance. Some ericoid mycorrhizal fungi have actually expanded their repertoire of enzymes for breaking down organic matter. They can extract nitrogen from cellulose, hemicellulose, lignin, pectin, and chitin. This would increase the benefit they can provide to their plant symbiotic partners.
Orchid mycorrhiza
All orchids are myco-heterotrophic at some stage during their lifecycle, meaning that they can survive only if they form orchid mycorrhizae. Orchid seeds are so small that they contain no nutrition to sustain the germinating seedling, and instead must gain the energy to grow from their fungal symbiont. The OM relationship is asymmetric; the plant seems to benefit more than the fungus, and some orchids are entirely mycoheterotrophic, lacking chlorophyll for photosynthesis. It is actually unknown whether fully autotrophic orchids that do not receive some of their carbon from fungi exist or not. Like fungi that form ErMs, OM fungi can sometimes live as endophytes or as independent saprotrophs. In the OM symbiosis, hyphae penetrate into the root cells and form pelotons (coils) for nutrient exchange.
Monotropoid mycorrhiza
This type of mycorrhiza occurs in the subfamily Monotropoideae of the Ericaceae, as well as several genera in the Orchidaceae. These plants are heterotrophic or mixotrophic and derive their carbon from the fungus partner. This is thus a non-mutualistic, parasitic type of mycorrhizal symbiosis.
Function
Mycorrhizal fungi form a mutualistic relationship with the roots of most plant species. In such a relationship, both the plants themselves and those parts of the roots that host the fungi, are said to be mycorrhizal. Relatively few of the mycorrhizal relationships between plant species and fungi have been examined to date, but 95% of the plant families investigated are predominantly mycorrhizal either in the sense that most of their species associate beneficially with mycorrhizae, or are absolutely dependent on mycorrhizae. The Orchidaceae are notorious as a family in which the absence of the correct mycorrhizae is fatal even to germinating seeds.
Recent research into ectomycorrhizal plants in boreal forests has indicated that mycorrhizal fungi and plants have a relationship that may be more complex than simply mutualistic. This relationship was noted when mycorrhizal fungi were unexpectedly found to be hoarding nitrogen from plant roots in times of nitrogen scarcity. Researchers argue that some mycorrhizae distribute nutrients based upon the environment with surrounding plants and other mycorrhizae. They go on to explain how this updated model could explain why mycorrhizae do not alleviate plant nitrogen limitation, and why plants can switch abruptly from a mixed strategy with both mycorrhizal and nonmycorrhizal roots to a purely mycorrhizal strategy as soil nitrogen availability declines. It has also been suggested that evolutionary and phylogenetic relationships can explain much more variation in the strength of mycorrhizal mutualisms than ecological factors.
Formation
To successfully engage in mutualistic symbiotic relationships with other organisms, such as mycorrhizal fungi and any of the thousands of microbes that colonize plants, plants must discriminate between mutualists and pathogens, allowing the mutualists to colonize while activating an immune response towards the pathogens. Plant genomes code for potentially hundreds of receptors for detecting chemical signals from other organisms. Plants dynamically adjust their symbiotic and immune responses, changing their interactions with their symbionts in response to feedbacks detected by the plant. In plants, the mycorrhizal symbiosis is regulated by the common symbiosis signaling pathway (CSSP), a set of genes involved in initiating and maintaining colonization by endosymbiotic fungi and other endosymbionts such as Rhizobia in legumes. The CSSP has origins predating the colonization of land by plants, demonstrating that the co-evolution of plants and arbuscular mycorrhizal fungi is over 500 million years old. In arbuscular mycorrhizal fungi, the presence of strigolactones, a plant hormone, secreted from roots induces fungal spores in the soil to germinate, stimulates their metabolism, growth and branching, and prompts the fungi to release chemical signals the plant can detect. Once the plant and fungus recognize one another as suitable symbionts, the plant activates the common symbiotic signaling pathway, which causes changes in the root tissues that enable the fungus to colonize.
Experiments with arbuscular mycorrhizal fungi have identified numerous chemical compounds to be involved in the "chemical dialog" that occurs between the prospective symbionts before symbiosis is begun. In plants, almost all plant hormones play a role in initiating or regulating AMF symbiosis, and other chemical compounds are also suspected to have a signaling function. While the signals emitted by the fungi are less understood, it has been shown that chitinaceous molecules known as Myc factors are essential for the formation of arbuscular mycorrhizae. Signals from plants are detected by LysM-containing receptor-like kinases, or LysM-RLKs. AMF genomes also code for potentially hundreds of effector proteins, of which only a few have a proven effect on mycorrhizal symbiosis, but many others likely have a function in communication with plant hosts as well.
Many factors are involved in the initiation of mycorrhizal symbiosis, but particularly influential is the plant's need for phosphorus. Experiments involving rice plants with a mutation disabling their ability to detect P starvation show that arbuscular mycorrhizal fungi detection, recruitment and colonization is prompted when the plant detects that it is starved of phosphorus. Nitrogen starvation also plays a role in initiating AMF symbiosis.
Mechanisms
The mechanisms by which mycorrhizae increase absorption include some that are physical and some that are chemical. Physically, most mycorrhizal mycelia are much smaller in diameter than the smallest root or root hair, and thus can explore soil material that roots and root hairs cannot reach, and provide a larger surface area for absorption. Chemically, the cell membrane chemistry of fungi differs from that of plants. For example, they may secrete organic acids that dissolve or chelate many ions, or release them from minerals by ion exchange. Mycorrhizae are especially beneficial for the plant partner in nutrient-poor soils.
Sugar-water/mineral exchange
The mycorrhizal mutualistic association provides the fungus with relatively constant and direct access to carbohydrates, such as glucose and sucrose. The carbohydrates are translocated from their source (usually leaves) to root tissue and on to the plant's fungal partners. In return, the plant gains the benefits of the mycelium's higher absorptive capacity for water and mineral nutrients, partly because of the large surface area of fungal hyphae, which are much longer and finer than plant root hairs, and partly because some such fungi can mobilize soil minerals unavailable to the plants' roots. The effect is thus to improve the plant's mineral absorption capabilities.
Unaided plant roots may be unable to take up nutrients that are chemically or physically immobilised; examples include phosphate ions and micronutrients such as iron. One form of such immobilization occurs in soil with high clay content, or soils with a strongly basic pH. The mycelium of the mycorrhizal fungus can, however, access many such nutrient sources, and make them available to the plants they colonize. Thus, many plants are able to obtain phosphate without using soil as a source. Another form of immobilisation is when nutrients are locked up in organic matter that is slow to decay, such as wood, and some mycorrhizal fungi act directly as decay organisms, mobilising the nutrients and passing some onto the host plants; for example, in some dystrophic forests, large amounts of phosphate and other nutrients are taken up by mycorrhizal hyphae acting directly on leaf litter, bypassing the need for soil uptake. Inga alley cropping, an agroforestry technique proposed as an alternative to slash and burn rainforest destruction, relies upon mycorrhiza within the root system of species of Inga to prevent the rain from washing phosphorus out of the soil.
In some more complex relationships, mycorrhizal fungi do not just collect immobilised soil nutrients, but connect individual plants together by mycorrhizal networks that transport water, carbon, and other nutrients directly from plant to plant through underground hyphal networks.
Suillus tomentosus, a basidiomycete fungus, produces specialized structures known as tuberculate ectomycorrhizae with its plant host lodgepole pine (Pinus contorta var. latifolia). These structures have been shown to host nitrogen fixing bacteria which contribute a significant amount of nitrogen and allow the pines to colonize nutrient-poor sites.
Disease, drought and salinity resistance and its correlation to mycorrhizae
Mycorrhizal plants are often more resistant to diseases, such as those caused by microbial soil-borne pathogens. These associations have been found to assist in plant defense both above and belowground. Mycorrhizas have been found to excrete enzymes that are toxic to soil borne organisms such as nematodes. More recent studies have shown that mycorrhizal associations result in a priming effect of plants that essentially acts as a primary immune response. When this association is formed a defense response is activated similarly to the response that occurs when the plant is under attack. As a result of this inoculation, defense responses are stronger in plants with mycorrhizal associations.
Ecosystem services provided by mycorrhizal fungi may depend on the soil microbiome. Furthermore, mycorrhizal fungi was significantly correlated with soil physical variable, but only with water level and not with aggregate stability and can lead also to more resistant to the effects of drought. Moreover, the significance of mycorrhizal fungi also includes alleviation of salt stress and its beneficial effects on plant growth and productivity. Although salinity can negatively affect mycorrhizal fungi, many reports show improved growth and performance of mycorrhizal plants under salt stress conditions.
Resistance to insects
Plants connected by mycorrhizal fungi in mycorrhizal networks can use these underground connections to communicate warning signals. For example, when a host plant is attacked by an aphid, the plant signals surrounding connected plants of its condition. Both the host plant and those connected to it release volatile organic compounds that repel aphids and attract parasitoid wasps, predators of aphids. This assists the mycorrhizal fungi by conserving its food supply.
Colonization of barren soil
Plants grown in sterile soils and growth media often perform poorly without the addition of spores or hyphae of mycorrhizal fungi to colonise the plant roots and aid in the uptake of soil mineral nutrients. The absence of mycorrhizal fungi can also slow plant growth in early succession or on degraded landscapes. The introduction of alien mycorrhizal plants to nutrient-deficient ecosystems puts indigenous non-mycorrhizal plants at a competitive disadvantage. This aptitude to colonize barren soil is defined by the category Oligotroph.
Resistance to toxicity
Fungi have a protective role for plants rooted in soils with high metal concentrations, such as acidic and contaminated soils. Pine trees inoculated with Pisolithus tinctorius planted in several contaminated sites displayed high tolerance to the prevailing contaminant, survivorship and growth. One study discovered the existence of Suillus luteus strains with varying tolerance of zinc. Another study discovered that zinc-tolerant strains of Suillus bovinus conferred resistance to plants of Pinus sylvestris. This was probably due to binding of the metal to the extramatricial mycelium of the fungus, without affecting the exchange of beneficial substances.
Occurrence of mycorrhizal associations
Mycorrhizas are present in 92% of plant families studied (80% of species), with arbuscular mycorrhizas being the ancestral and predominant form, and the most prevalent symbiotic association found in the plant kingdom. The structure of arbuscular mycorrhizas has been highly conserved since their first appearance in the fossil record, with both the development of ectomycorrhizas and the loss of mycorrhizas, evolving convergently on multiple occasions.
Associations of fungi with the roots of plants have been known since at least the mid-19th century. However, early observers simply recorded the fact without investigating the relationships between the two organisms. This symbiosis was studied and described by Franciszek Kamieński in 1879–1882.
Climate change
CO2 released by human activities is causing climate change and possible damage to mycorrhizae, but the direct effect of an increase in the gas should be to benefit plants and mycorrhizae. In Arctic regions, nitrogen and water are harder for plants to obtain, making mycorrhizae crucial to plant growth. Since mycorrhizae tend to do better in cooler temperatures, warming could be detrimental to them. Gases such as SO2, NO-x, and O3 produced by human activity may harm mycorrhizae, causing reduction in "propagules, the colonization of roots, degradation in connections between trees, reduction in the mycorrhizal incidence in trees, and reduction in the enzyme activity of ectomycorrhizal roots."
A company in Israel, Groundwork BioAg, has discovered a method of using mycorrhizal fungi to increase agricultural crops while sequestering greenhouse gases and eliminating CO2 from the atmosphere.
Conservation and mapping
In 2021, the Society for the Protection of Underground Networks was launched. SPUN is a science-based initiative to map and protect the mycorrhizal networks regulating Earth’s climate and ecosystems. Its stated goals are mapping, protecting, and harnessing mycorrhizal fungi.
| Biology and health sciences | Basics | Plants |
59364 | https://en.wikipedia.org/wiki/Methionine | Methionine | Methionine (symbol Met or M) () is an essential amino acid in humans.
As the precursor of other non-essential amino acids such as cysteine and taurine, versatile compounds such as SAM-e, and the important antioxidant glutathione, methionine plays a critical role in the metabolism and health of many species, including humans. Methionine is also involved in angiogenesis and various processes related to DNA transcription, epigenetic expression, and gene regulation.
Methionine was first isolated in 1921 by John Howard Mueller. It is encoded by the codon AUG. It was named by Satoru Odake in 1925, as an abbreviation of its structural description 2-amino-4-(methylthio)butanoic acid.
Biochemical details
Methionine (abbreviated as Met or M; encoded by the codon AUG) is an α-amino acid that is used in the biosynthesis of proteins. It contains a carboxyl group (which is in the deprotonated −COO− form under biological pH conditions), an amino group (which is in the protonated form under biological pH conditions) located in α-position with respect to the carboxyl group, and an S-methyl thioether side chain, classifying it as a nonpolar, aliphatic amino acid.
In nuclear genes of eukaryotes and in Archaea, methionine is coded for by the start codon, meaning it indicates the start of the coding region and is the first amino acid produced in a nascent polypeptide during mRNA translation.
A proteinogenic amino acid
Cysteine and methionine are the two sulfur-containing proteinogenic amino acids. Excluding the few exceptions where methionine may act as a redox sensor (e.g.,methionine sulfoxide), methionine residues do not have a catalytic role. This is in contrast to cysteine residues, where the thiol group has a catalytic role in many proteins. The thioether within methionine does however have a minor structural role due to the stability effect of S/π interactions between the side chain sulfur atom and aromatic amino acids in one-third of all known protein structures. This lack of a strong role is reflected in experiments where little effect is seen in proteins where methionine is replaced by norleucine, a straight hydrocarbon sidechain amino acid which lacks the thioether.
It has been conjectured that norleucine was present in early versions of the genetic code, but methionine intruded into the final version of the genetic code due to the fact it is used in the cofactor S-adenosylmethionine (SAM-e). This situation is not unique and may have occurred with ornithine and arginine.
Encoding
Methionine is one of only two amino acids encoded by a single codon (AUG) in the standard genetic code (tryptophan, encoded by UGG, is the other). In reflection to the evolutionary origin of its codon, the other AUN codons encode isoleucine, which is also a hydrophobic amino acid. In the mitochondrial genome of several organisms, including metazoa and yeast, the codon AUA also encodes for methionine. In the standard genetic code AUA codes for isoleucine and the respective tRNA (ileX in Escherichia coli) uses the unusual base lysidine (bacteria) or agmatidine (archaea) to discriminate against AUG.
The methionine codon AUG is also the most common start codon. A "Start" codon is message for a ribosome that signals the initiation of protein translation from mRNA when the AUG codon is in a Kozak consensus sequence. As a consequence, methionine is often incorporated into the N-terminal position of proteins in eukaryotes and archaea during translation, although it can be removed by post-translational modification. In bacteria, the derivative N-formylmethionine is used as the initial amino acid.
Derivatives
S-Adenosylmethionine
The methionine-derivative S-adenosylmethionine (SAM-e) is a cofactor that serves mainly as a methyl donor. SAM-e is composed of an adenosyl molecule (via 5′ carbon) attached to the sulfur of methionine, therefore making it a sulfonium cation (i.e., three substituents and positive charge). The sulfur acts as a soft Lewis acid (i.e., donor/electrophile) which allows the S-methyl group to be transferred to an oxygen, nitrogen, or aromatic system, often with the aid of other cofactors such as cobalamin (vitamin B12 in humans). Some enzymes use SAM-e to initiate a radical reaction; these are called radical SAM-e enzymes.
As a result of the transfer of the methyl group, S-adenosylhomocysteine is obtained. In bacteria, this is either regenerated by methylation or is salvaged by removing the adenine and the homocysteine, leaving the compound dihydroxypentandione to spontaneously convert into autoinducer-2, which is excreted as a waste product or quorum signal.
Biosynthesis
As an essential amino acid, methionine is not synthesized de novo in humans and other animals, which must ingest methionine or methionine-containing proteins. In plants and microorganisms, methionine biosynthesis belongs to the aspartate family, along with threonine and lysine (via diaminopimelate, but not via α-aminoadipate). The main backbone is derived from aspartic acid, while the sulfur may come from cysteine, methanethiol, or hydrogen sulfide.
First, aspartic acid is converted via β-aspartyl semialdehyde into homoserine by two reduction steps of the terminal carboxyl group (homoserine has therefore a γ-hydroxyl, hence the homo- series). The intermediate aspartate semialdehyde is the branching point with the lysine biosynthetic pathway, where it is instead condensed with pyruvate. Homoserine is the branching point with the threonine pathway, where instead it is isomerised after activating the terminal hydroxyl with phosphate (also used for methionine biosynthesis in plants).
Homoserine is then activated with a phosphate, succinyl or an acetyl group on the hydroxyl.
In plants and possibly in some bacteria, phosphate is used. This step is shared with threonine biosynthesis.
In most organisms, an acetyl group is used to activate the homoserine. This can be catalysed in bacteria by an enzyme encoded by metX or metA (not homologues).
In enterobacteria and a limited number of other organisms, succinate is used. The enzyme that catalyses the reaction is MetA and the specificity for acetyl-CoA and succinyl-CoA is dictated by a single residue. The physiological basis for the preference of acetyl-CoA or succinyl-CoA is unknown, but such alternative routes are present in some other pathways (e.g. lysine biosynthesis and arginine biosynthesis).
The hydroxyl activating group is then replaced with cysteine, methanethiol, or hydrogen sulfide. A replacement reaction is technically a γ-elimination followed by a variant of a Michael addition. All the enzymes involved are homologues and members of the Cys/Met metabolism PLP-dependent enzyme family, which is a subset of the PLP-dependent fold type I clade. They utilise the cofactor PLP (pyridoxal phosphate), which functions by stabilising carbanion intermediates.
If it reacts with cysteine, it produces cystathionine, which is cleaved to yield homocysteine. The enzymes involved are cystathionine-γ-synthase (encoded by metB in bacteria) and cystathionine-β-lyase (metC). Cystathionine is bound differently in the two enzymes allowing β or γ reactions to occur.
If it reacts with free hydrogen sulfide, it produces homocysteine. This is catalysed by O-acetylhomoserine aminocarboxypropyltransferase (formerly known as O-acetylhomoserine (thiol)-lyase. It is encoded by either metY or metZ in bacteria.
If it reacts with methanethiol, it produces methionine directly. Methanethiol is a byproduct of catabolic pathway of certain compounds, therefore this route is more uncommon.
If homocysteine is produced, the thiol group is methylated, yielding methionine. Two methionine synthases are known; one is cobalamin (vitamin B12) dependent and one is independent.
The pathway using cysteine is called the "transsulfuration pathway", while the pathway using hydrogen sulfide (or methanethiol) is called "direct-sulfurylation pathway".
Cysteine is similarly produced, namely it can be made from an activated serine and either from homocysteine ("reverse transsulfurylation route") or from hydrogen sulfide ("direct sulfurylation route"); the activated serine is generally O-acetylserine (via CysK or CysM in E. coli), but in Aeropyrum pernix and some other archaea O-phosphoserine is used. CysK and CysM are homologues, but belong to the PLP fold type III clade.
Transsulfurylation pathway
Enzymes involved in the E. coli transsulfurylation route of methionine biosynthesis:
Aspartokinase
Aspartate-semialdehyde dehydrogenase
Homoserine dehydrogenase
Homoserine O-transsuccinylase
Cystathionine-γ-synthase
Cystathionine-β-lyase
Methionine synthase (in mammals, this step is performed by homocysteine methyltransferase or betaine—homocysteine S-methyltransferase.)
Other biochemical pathways
Although mammals cannot synthesize methionine, they can still use it in a variety of biochemical pathways:
Catabolism
Methionine is converted to S-adenosylmethionine (SAM-e) by (1) methionine adenosyltransferase.
SAM-e serves as a methyl donor in many (2) methyltransferase reactions, and is converted to S-adenosylhomocysteine (SAH).
(3) Adenosylhomocysteinase
cysteine.
Regeneration
Methionine can be regenerated from homocysteine via (4) methionine synthase in a reaction that requires vitamin B12 as a cofactor.
Homocysteine can also be remethylated using glycine betaine (N,N,N-trimethylglycine, TMG) to methionine via the enzyme betaine-homocysteine methyltransferase (E.C.2.1.1.5, BHMT). BHMT makes up to 1.5% of all the soluble protein of the liver, and recent evidence suggests that it may have a greater influence on methionine and homocysteine homeostasis than methionine synthase.
Reverse-transulfurylation pathway: conversion to cysteine
Homocysteine can be converted to cysteine.
(5) Cystathionine-β-synthase (an enzyme which requires pyridoxal phosphate, the active form of vitamin B6) combines homocysteine and serine to produce cystathionine. Instead of degrading cystathionine via cystathionine-β-lyase, as in the biosynthetic pathway, cystathionine is broken down to cysteine and α-ketobutyrate via (6) cystathionine-γ-lyase.
(7) The enzyme α-ketoacid dehydrogenase converts α-ketobutyrate to propionyl-CoA, which is metabolized to succinyl-CoA in a three-step process (see propionyl-CoA for pathway).
Ethylene synthesis
This amino acid is also used by plants for synthesis of ethylene. The process is known as the Yang cycle or the methionine cycle.
Metabolic diseases
The degradation of methionine is impaired in the following metabolic diseases:
Combined malonic and methylmalonic aciduria (CMAMMA)
Homocystinuria
Methylmalonic acidemia
Propionic acidemia
Chemical synthesis
The industrial synthesis combines acrolein, methanethiol, and cyanide, which affords the hydantoin. Racemic methionine can also be synthesized from diethyl sodium phthalimidomalonate by alkylation with chloroethylmethylsulfide (ClCH2CH2SCH3) followed by hydrolysis and decarboxylation. Also see Methanol.
Human nutrition
There is inconclusive clinical evidence on methionine supplementation. Dietary restriction of methionine can lead to bone-related disorders.
Methionine supplementation may benefit those suffering from copper poisoning.
Overconsumption of methionine, the methyl group donor in DNA methylation, is related to cancer growth in a number of studies.
Requirements
The Food and Nutrition Board of the U.S. Institute of Medicine set Recommended Dietary Allowances (RDAs) for essential amino acids in 2002. For methionine combined with cysteine, for adults 19 years and older, 19 mg/kg body weight/day.
This translates to about 1.33 grams per day for a 70 kilogram individual.
Dietary sources
High levels of methionine can be found in eggs, meat, and fish; sesame seeds, Brazil nuts, and some other plant seeds; and cereal grains. Most fruits and vegetables contain very little. Most legumes, though protein dense, are low in methionine. Proteins without adequate methionine are not considered to be complete proteins. For that reason, racemic methionine is sometimes added as an ingredient to pet foods.
Health
Loss of methionine has been linked to senile greying of hair. Its lack leads to a buildup of hydrogen peroxide in hair follicles, a reduction in tyrosinase effectiveness, and a gradual loss of hair color. Methionine raises the intracellular concentration of glutathione, thereby promoting antioxidant-mediated cell defense and redox regulation. It also protects cells against dopamine induced nigral cell loss by binding oxidative metabolites.
Methionine is an intermediate in the biosynthesis of cysteine, carnitine, taurine, lecithin, phosphatidylcholine, and other phospholipids. Improper conversion of methionine can lead to atherosclerosis due to accumulation of homocysteine.
Other uses
DL-Methionine is sometimes given as a supplement to dogs; It helps reduce the chances of kidney stones in dogs. Methionine is also known to increase the urinary excretion of quinidine by acidifying the urine. Aminoglycoside antibiotics used to treat urinary tract infections work best in alkaline conditions, and urinary acidification from using methionine can reduce its effectiveness. If a dog is on a diet that acidifies the urine, methionine should not be used.
Methionine is allowed as a supplement to organic poultry feed under the US certified organic program.
Methionine can be used as a nontoxic pesticide option against giant swallowtail caterpillars, which are a serious pest to orange crops.
| Biology and health sciences | Amino acids | Biology |
59366 | https://en.wikipedia.org/wiki/Large%20intestine | Large intestine | The large intestine, also known as the large bowel, is the last part of the gastrointestinal tract and of the digestive system in tetrapods. Water is absorbed here and the remaining waste material is stored in the rectum as feces before being removed by defecation. The colon (progressing from the ascending colon to the transverse, the descending and finally the sigmoid colon) is the longest portion of the large intestine, and the terms "large intestine" and "colon" are often used interchangeably, but most sources define the large intestine as the combination of the cecum, colon, rectum, and anal canal. Some other sources exclude the anal canal.
In humans, the large intestine begins in the right iliac region of the pelvis, just at or below the waist, where it is joined to the end of the small intestine at the cecum, via the ileocecal valve. It then continues as the colon ascending the abdomen, across the width of the abdominal cavity as the transverse colon, and then descending to the rectum and its endpoint at the anal canal. Overall, in humans, the large intestine is about long, which is about one-fifth of the whole length of the human gastrointestinal tract.
Structure
The colon of the large intestine is the last part of the digestive system. It has a segmented appearance due to a series of saccules called haustra. It extracts water and salt from solid wastes before they are eliminated from the body and is the site in which the fermentation of unabsorbed material by the gut microbiota occurs. Unlike the small intestine, the colon does not play a major role in absorption of foods and nutrients. About 1.5 litres or 45 ounces of water arrives in the colon each day.
The colon is the longest part of the large intestine and its average length in the adult human is 65 inches or 166 cm (range of 80 to 313 cm) for males, and 61 inches or 155 cm (range of 80 to 214 cm) for females.
Sections
In mammals, the large intestine consists of the cecum (including the appendix), colon (the longest part), rectum, and anal canal.
The four sections of the colon are: the ascending colon, transverse colon, descending colon, and sigmoid colon. These sections turn at the colic flexures.
The parts of the colon are either intraperitoneal or behind it in the retroperitoneum. Retroperitoneal organs, in general, do not have a complete covering of peritoneum, so they are fixed in location. Intraperitoneal organs are completely surrounded by peritoneum and are therefore mobile. Of the colon, the ascending colon, descending colon and rectum are retroperitoneal, while the cecum, appendix, transverse colon and sigmoid colon are intraperitoneal. This is important as it affects which organs can be easily accessed during surgery, such as a laparotomy.
In terms of diameter, the cecum is the widest, averaging slightly less than 9 cm in healthy individuals, and the transverse colon averages less than 6 cm in diameter. The descending and sigmoid colon are slightly smaller, with the sigmoid colon averaging in diameter. Diameters larger than certain thresholds for each colonic section can be diagnostic for megacolon.
Cecum and appendix
The cecum is the first section of the large intestine and is involved in digestion, while the appendix which develops embryologically from it, is not involved in digestion and is considered to be part of the gut-associated lymphoid tissue. The function of the appendix is uncertain, but some sources believe that it has a role in housing a sample of the gut microbiota, and is able to help to repopulate the colon with microbiota if depleted during the course of an immune reaction. The appendix has also been shown to have a high concentration of lymphatic cells.
Ascending colon
The ascending colon is the first of four main sections of the large intestine. It is connected to the small intestine by a section of bowel called the cecum. The ascending colon runs upwards through the abdominal cavity toward the transverse colon for approximately eight inches (20 cm).
One of the main functions of the colon is to remove the water and other key nutrients from waste material and recycle it. As the waste material exits the small intestine through the ileocecal valve, it will move into the cecum and then to the ascending colon where this process of extraction starts. The waste material is pumped upwards toward the transverse colon by peristalsis. The ascending colon is sometimes attached to the appendix via Gerlach's valve. In ruminants, the ascending colon is known as the spiral colon.
Taking into account all ages and sexes, colon cancer occurs here most often (41%).
Transverse colon
The transverse colon is the part of the colon from the hepatic flexure, also known as the right colic, (the turn of the colon by the liver) to the splenic flexure also known as the left colic, (the turn of the colon by the spleen). The transverse colon hangs off the stomach, attached to it by a large fold of peritoneum called the greater omentum. On the posterior side, the transverse colon is connected to the posterior abdominal wall by a mesentery known as the transverse mesocolon.
The transverse colon is encased in peritoneum, and is therefore mobile (unlike the parts of the colon immediately before and after it).
The proximal two-thirds of the transverse colon is perfused by the middle colic artery, a branch of the superior mesenteric artery (SMA), while the latter third is supplied by branches of the inferior mesenteric artery (IMA). The "watershed" area between these two blood supplies, which represents the embryologic division between the midgut and hindgut, is an area sensitive to ischemia.
Descending colon
The descending colon is the part of the colon from the splenic flexure to the beginning of the sigmoid colon. One function of the descending colon in the digestive system is to store feces that will be emptied into the rectum. It is retroperitoneal in two-thirds of humans. In the other third, it has a (usually short) mesentery. The arterial supply comes via the left colic artery. The descending colon is also called the distal gut, as it is further along the gastrointestinal tract than the proximal gut. Gut flora are very dense in this region.
Sigmoid colon
The sigmoid colon is the part of the large intestine after the descending colon and before the rectum. The name sigmoid means S-shaped (see sigmoid; cf. sigmoid sinus). The walls of the sigmoid colon are muscular and contract to increase the pressure inside the colon, causing the stool to move into the rectum.
The sigmoid colon is supplied with blood from several branches (usually between 2 and 6) of the sigmoid arteries, a branch of the IMA. The IMA terminates as the superior rectal artery.
Sigmoidoscopy is a common diagnostic technique used to examine the sigmoid colon.
Rectum
The rectum is the last section of the large intestine. It holds the formed feces awaiting elimination via defecation.
It is about 12 cm long.
Appearance
The cecum – the first part of the large intestine
Taeniae coli – three bands of smooth muscle
Haustra – bulges caused by contraction of taeniae coli
Epiploic appendages – small fat accumulations on the viscera
The taenia coli run the length of the large intestine. Because the taenia coli are shorter than the large bowel itself, the colon becomes sacculated, forming the haustra of the colon which are the shelf-like intraluminal projections.
Blood supply
Arterial supply to the colon comes from branches of the superior mesenteric artery (SMA) and inferior mesenteric artery (IMA). Flow between these two systems communicates via the marginal artery of the colon that runs parallel to the colon for its entire length. Historically, a structure variously identified as the arc of Riolan or meandering mesenteric artery (of Moskowitz) was thought to connect the proximal SMA to the proximal IMA. This variably present structure would be important if either vessel were occluded. However, at least one review of the literature questions the existence of this vessel, with some experts calling for the abolition of these terms from future medical literature.
Venous drainage usually mirrors colonic arterial supply, with the inferior mesenteric vein draining into the splenic vein, and the superior mesenteric vein joining the splenic vein to form the hepatic portal vein that then enters the liver. Middle rectal veins are an exception, delivering blood to inferior vena cava and bypassing the liver.
Lymphatic drainage
Lymphatic drainage from the ascending colon and proximal two-thirds of the transverse colon is to the ileocolic lymph nodes and the superior mesenteric lymph nodes, which drain into the cisterna chyli. The lymph from the distal one-third of the transverse colon, the descending colon, the sigmoid colon, and the upper rectum drain into the inferior mesenteric and colic lymph nodes. The lower rectum to the anal canal above the pectinate line drain to the internal ileocolic nodes. The anal canal below the pectinate line drains into the superficial inguinal nodes. The pectinate line only roughly marks this transition.
Nerve supply
Sympathetic supply: superior & inferior mesenteric ganglia;
parasympathetic supply: vagus & sacral plexus (S2-S4)
Development
The endoderm, mesoderm and ectoderm are germ layers that develop in a process called gastrulation. Gastrulation occurs early in human development. The gastrointestinal tract is derived from these layers.
Variation
One variation on the normal anatomy of the colon occurs when extra loops form, resulting in a colon that is up to five metres longer than normal. This condition, referred to as redundant colon, typically has no direct major health consequences, though rarely volvulus occurs, resulting in obstruction and requiring immediate medical attention. A significant indirect health consequence is that use of a standard adult colonoscope is difficult and in some cases impossible when a redundant colon is present, though specialized variants on the instrument (including the pediatric variant) are useful in overcoming this problem.
Microanatomy
Colonic crypts
The wall of the large intestine is lined with simple columnar epithelium with invaginations. The invaginations are called the intestinal glands or colonic crypts.
The colon crypts are shaped like microscopic thick walled test tubes with a central hole down the length of the tube (the crypt lumen). Four tissue sections are shown here, two cut across the long axes of the crypts and two cut parallel to the long axes. In these images the cells have been stained by immunohistochemistry to show a brown-orange color if the cells produce a mitochondrial protein called cytochrome c oxidase subunit I (CCOI). The nuclei of the cells (located at the outer edges of the cells lining the walls of the crypts) are stained blue-gray with haematoxylin. As seen in panels C and D, crypts are about 75 to about 110 cells long. Baker et al. found that the average crypt circumference is 23 cells. Thus, by the images shown here, there are an average of about 1,725 to 2,530 cells per colonic crypt. Nooteboom et al. measuring the number of cells in a small number of crypts reported a range of 1,500 to 4,900 cells per colonic crypt. Cells are produced at the crypt base and migrate upward along the crypt axis before being shed into the colonic lumen days later. There are 5 to 6 stem cells at the bases of the crypts.
As estimated from the image in panel A, there are about 100 colonic crypts per square millimeter of the colonic epithelium. Since the average length of the human colon is 160.5 cm and the average inner circumference of the colon is 6.2 cm, the inner surface epithelial area of the human colon has an average area of about 995 cm2, which includes 9,950,000 (close to 10 million) crypts.
In the four tissue sections shown here, many of the intestinal glands have cells with a mitochondrial DNA mutation in the CCOI gene and appear mostly white, with their main color being the blue-gray staining of the nuclei. As seen in panel B, a portion of the stem cells of three crypts appear to have a mutation in CCOI, so that 40% to 50% of the cells arising from those stem cells form a white segment in the cross cut area.
Overall, the percent of crypts deficient for CCOI is less than 1% before age 40, but then increases linearly with age. Colonic crypts deficient for CCOI in women reaches, on average, 18% in women and 23% in men by 80–84 years of age.
Crypts of the colon can reproduce by fission, as seen in panel C, where a crypt is fissioning to form two crypts, and in panel B where at least one crypt appears to be fissioning. Most crypts deficient in CCOI are in clusters of crypts (clones of crypts) with two or more CCOI-deficient crypts adjacent to each other (see panel D).
Mucosa
About 150 of the many thousands of protein coding genes expressed in the large intestine, some are specific to the mucous membrane in different regions and include CEACAM7.
Function
The large intestine absorbs water and any remaining absorbable nutrients from the food before sending the indigestible matter to the rectum. The colon absorbs vitamins that are created by the colonic bacteria, such as thiamine, riboflavin, and vitamin K (especially important as the daily ingestion of vitamin K is not normally enough to maintain adequate blood coagulation). It also compacts feces, and stores fecal matter in the rectum until it can be discharged via the anus in defecation.
The large intestine also secretes K+ and Cl-. Chloride secretion increases in cystic fibrosis.
Recycling of various nutrients takes place in the colon. Examples include fermentation of carbohydrates, short chain fatty acids, and urea cycling.
The appendix contains a small amount of mucosa-associated lymphoid tissue which gives the appendix an undetermined role in immunity. However, the appendix is known to be important in fetal life as it contains endocrine cells that release biogenic amines and peptide hormones important for homeostasis during early growth and development.
By the time the chyme has reached this tube, most nutrients and 90% of the water have been absorbed by the body. Indeed, as demonstrated by the commonality of ileostomy procedures, it is possible for many people to live without large portions of their large intestine, or even without it completely. At this point only some electrolytes like sodium, magnesium, and chloride are left as well as indigestible parts of ingested food (e.g., a large part of ingested amylose, starch which has been shielded from digestion heretofore, and dietary fiber, which is largely indigestible carbohydrate in either soluble or insoluble form). As the chyme moves through the large intestine, most of the remaining water is removed, while the chyme is mixed with mucus and bacteria (known as gut flora), and becomes feces. The ascending colon receives fecal material as a liquid. The muscles of the colon then move the watery waste material forward and slowly absorb all the excess water, causing the stools to gradually solidify as they move along into the descending colon.
The bacteria break down some of the fiber for their own nourishment and create acetate, propionate, and butyrate as waste products, which in turn are used by the cell lining of the colon for nourishment. No protein is made available. In humans, perhaps 10% of the undigested carbohydrate thus becomes available, though this may vary with diet; in other animals, including other apes and primates, who have proportionally larger colons, more is made available, thus permitting a higher portion of plant material in the diet. The large intestine produces no digestive enzymes — chemical digestion is completed in the small intestine before the chyme reaches the large intestine. The pH in the colon varies between 5.5 and 7 (slightly acidic to neutral).
Standing gradient osmosis
Water absorption at the colon typically proceeds against a transmucosal osmotic pressure gradient. The standing gradient osmosis is the reabsorption of water against the osmotic gradient in the intestines. Cells occupying the intestinal lining pump sodium ions into the intercellular space, raising the osmolarity of the intercellular fluid. This hypertonic fluid creates an osmotic pressure that drives water into the lateral intercellular spaces by osmosis via tight junctions and adjacent cells, which then in turn moves across the basement membrane and into the capillaries, while more sodium ions are pumped again into the intercellular fluid. Although water travels down an osmotic gradient in each individual step, overall, water usually travels against the osmotic gradient due to the pumping of sodium ions into the intercellular fluid. This allows the large intestine to absorb water despite the blood in capillaries being hypotonic compared to the fluid within the intestinal lumen.
Gut flora
The large intestine houses over 700 species of bacteria that perform a variety of functions, as well as fungi, protozoa, and archaea. Species diversity varies by geography and diet. The microbes in a human distal gut often number in the vicinity of 100 trillion, and can weigh around 200 grams (0.44 pounds). This mass of mostly symbiotic microbes has recently been called the latest human organ to be "discovered" or in other words, the "forgotten organ".
The large intestine absorbs some of the products formed by the bacteria inhabiting this region. Undigested polysaccharides (fiber) are metabolized to short-chain fatty acids by bacteria in the large intestine and absorbed by passive diffusion. The bicarbonate that the large intestine secretes helps to neutralize the increased acidity resulting from the formation of these fatty acids.
These bacteria also produce large amounts of vitamins, especially vitamin K and biotin (a B vitamin), for absorption into the blood. Although this source of vitamins, in general, provides only a small part of the daily requirement, it makes a significant contribution when dietary vitamin intake is low. An individual who depends on absorption of vitamins formed by bacteria in the large intestine may become vitamin-deficient if treated with antibiotics that inhibit the vitamin producing species of bacteria as well as the intended disease-causing bacteria.
Other bacterial products include gas (flatus), which is a mixture of nitrogen and carbon dioxide, with small amounts of the gases hydrogen, methane, and hydrogen sulfide. Bacterial fermentation of undigested polysaccharides produces these. Some of the fecal odor is due to indoles, metabolized from the amino acid tryptophan. The normal flora is also essential in the development of certain tissues, including the cecum and lymphatics.
They are also involved in the production of cross-reactive antibodies. These are antibodies produced by the immune system against the normal flora, that are also effective against related pathogens, thereby preventing infection or invasion.
The two most prevalent phyla of the colon are Bacillota and Bacteroidota. The ratio between the two seems to vary widely as reported by the Human Microbiome Project. Bacteroides are implicated in the initiation of colitis and colon cancer. Bifidobacteria are also abundant, and are often described as 'friendly bacteria'.
A mucus layer protects the large intestine from attacks from colonic commensal bacteria.
Clinical significance
Disease
Following are the most common diseases or disorders of the colon:
Colonoscopy
Colonoscopy is the endoscopic examination of the large intestine and the distal part of the small bowel with a CCD camera or a fiber optic camera on a flexible tube passed through the anus. It can provide a visual diagnosis (e.g. ulceration, polyps) and grants the opportunity for biopsy or removal of suspected colorectal cancer lesions. Colonoscopy can remove polyps as small as one millimetre or less. Once polyps are removed, they can be studied with the aid of a microscope to determine if they are precancerous or not. It takes 15 years or fewer for a polyp to turn cancerous.
Colonoscopy is similar to sigmoidoscopy—the difference being related to which parts of the colon each can examine. A colonoscopy allows an examination of the entire colon (1200–1500 mm in length). A sigmoidoscopy allows an examination of the distal portion (about 600 mm) of the colon, which may be sufficient because benefits to cancer survival of colonoscopy have been limited to the detection of lesions in the distal portion of the colon.
A sigmoidoscopy is often used as a screening procedure for a full colonoscopy, often done in conjunction with a stool-based test such as a fecal occult blood test (FOBT), fecal immunochemical test (FIT), or multi-target stool DNA test (Cologuard) or blood-based test, SEPT9 DNA methylation test (Epi proColon). About 5% of these screened patients are referred to colonoscopy.
Virtual colonoscopy, which uses 2D and 3D imagery reconstructed from computed tomography (CT) scans or from nuclear magnetic resonance (MR) scans, is also possible, as a totally non-invasive medical test, although it is not standard and still under investigation regarding its diagnostic abilities. Furthermore, virtual colonoscopy does not allow for therapeutic maneuvers such as polyp/tumour removal or biopsy nor visualization of lesions smaller than 5 millimeters. If a growth or polyp is detected using CT colonography, a standard colonoscopy would still need to be performed. Additionally, surgeons have lately been using the term pouchoscopy to refer to a colonoscopy of the ileo-anal pouch.
Other animals
The large intestine is truly distinct only in tetrapods, in which it is almost always separated from the small intestine by an ileocaecal valve. In most vertebrates, however, it is a relatively short structure running directly to the anus, although noticeably wider than the small intestine. Although the caecum is present in most amniotes, only in mammals does the remainder of the large intestine develop into a true colon.
In some small mammals, the colon is straight, as it is in other tetrapods, but, in the majority of mammalian species, it is divided into ascending and descending portions; a distinct transverse colon is typically present only in primates. However, the taeniae coli and accompanying haustra are not found in either carnivorans or ruminants. The rectum of mammals (other than monotremes) is derived from the cloaca of other vertebrates, and is, therefore, not truly homologous with the "rectum" found in these species.
In some fish, there is no true large intestine, but simply a short rectum connecting the end of the digestive part of the gut to the cloaca. In sharks, this includes a rectal gland that secretes salt to help the animal maintain osmotic balance with the seawater. The gland somewhat resembles a caecum in structure but is not a homologous structure.
Additional images
| Biology and health sciences | Digestive system | null |
59385 | https://en.wikipedia.org/wiki/Amanita%20muscaria | Amanita muscaria | Amanita muscaria, commonly known as the fly agaric or fly amanita, is a basidiomycete of the genus Amanita. It is a large white-gilled, white-spotted, and usually red mushroom.
Despite its easily distinguishable features, A.muscaria is a fungus with several known variations, or subspecies. These subspecies are slightly different, some having yellow or white caps, but are all usually called fly agarics, most often recognizable by their notable white spots. Recent DNA fungi research, however, has shown that some mushrooms called "fly agaric" are in fact unique species, such as A.persicina (the peach-colored fly agaric).
Native throughout the temperate and boreal regions of the Northern Hemisphere, A.muscaria has been unintentionally introduced to many countries in the Southern Hemisphere, generally as a symbiont with pine and birch plantations, and is now a true cosmopolitan species. It associates with various deciduous and coniferous trees.
Although poisonous, death due to poisoning from A.muscaria ingestion is quite rare. Parboiling twice with water weakens its toxicity and breaks down the mushroom's psychoactive substances; it is eaten in parts of Europe, Asia, and North America. All A.muscaria varieties, but in particular A.muscaria var. muscaria, are noted for their hallucinogenic properties, with the main psychoactive constituents being muscimol and its neurotoxic precursor ibotenic acid. A local variety of the mushroom was used as an intoxicant and entheogen by the indigenous peoples of Siberia.
Arguably the most iconic toadstool species, the fly agaric is one of the most recognizable fungi in the world, and is widely encountered in popular culture, including in video games—for example, the frequent use of a recognizable A.muscaria in the Mario franchise (e.g. its Super Mushroom power-up)—and television—for example, the houses in The Smurfs franchise. There have been cases of children admitted to hospitals after consuming this poisonous mushroom; the children may have been attracted to it because of its pop-culture associations.
Taxonomy
The name of the mushroom in many European languages is thought to derive from its use as an insecticide when sprinkled in milk. This practice has been recorded from Germanic- and Slavic-speaking parts of Europe, as well as the Vosges region and pockets elsewhere in France, and Romania. Albertus Magnus was the first to record it in his work De vegetabilibus some time before 1256, commenting "vocatur fungus muscarum, eo quod in lacte pulverizatus interficit muscas" ("it is called the fly mushroom because it is powdered in milk to kill flies").
The 16th-century Flemish botanist Carolus Clusius traced the practice of sprinkling it into milk to Frankfurt in Germany, while Carl Linnaeus, the "father of taxonomy", reported it from Småland in southern Sweden, where he had lived as a child. He described it in volume two of his Species Plantarum in 1753, giving it the name Agaricus muscarius, the specific epithet deriving from Latin musca meaning "fly". It gained its current name in 1783, when placed in the genus Amanita by Jean-Baptiste Lamarck, a name sanctioned in 1821 by the "father of mycology", Swedish naturalist Elias Magnus Fries. The starting date for all the mycota had been set by general agreement as January 1, 1821, the date of Fries's work, and so the full name was then Amanita muscaria (L.:Fr.) Hook. The 1987 edition of the International Code of Botanical Nomenclature changed the rules on the starting date and primary work for names of fungi, and names can now be considered valid as far back as May 1, 1753, the date of publication of Linnaeus's work. Hence, Linnaeus and Lamarck are now taken as the namers of Amanita muscaria (L.) Lam..
The English mycologist John Ramsbottom reported that Amanita muscaria was used for getting rid of bugs in England and Sweden, and bug agaric was an old alternative name for the species. French mycologist Pierre Bulliard reported having tried without success to replicate its fly-killing properties in his work (1784), and proposed a new binomial name Agaricus pseudo-aurantiacus because of this. One compound isolated from the fungus is 1,3-diolein (1,3-di(cis-9-octadecenoyl)glycerol), which attracts insects.
It has been hypothesised that the flies intentionally seek out the fly agaric for its intoxicating properties.
An alternative derivation proposes that the term fly- refers not to insects as such but rather the delirium resulting from consumption of the fungus. This is based on the medieval belief that flies could enter a person's head and cause mental illness. Several regional names appear to be linked with this connotation, meaning the "mad" or "fool's" version of the highly regarded edible mushroom Amanita caesarea. Hence there is "mad oriol" in Catalan, mujolo folo from Toulouse, from the Aveyron department in Southern France, from Trentino in Italy. A local dialect name in Fribourg in Switzerland is tsapi de diablhou, which translates as "Devil's hat".
Classification
Amanita muscaria is the type species of the genus. By extension, it is also the type species of Amanita subgenus Amanita, as well as section Amanita within this subgenus. Amanita subgenus Amanita includes all Amanita with inamyloid spores. Amanita section Amanita includes the species with patchy universal veil remnants, including a volva that is reduced to a series of concentric rings, and the veil remnants on the cap to a series of patches or warts. Most species in this group also have a bulbous base. Amanita section Amanita consists of A. muscaria and its close relatives, including A. pantherina (the panther cap), A. gemmata, A. farinosa, and A. xanthocephala. Modern fungal taxonomists have classified Amanita muscaria and its allies this way based on gross morphology and spore inamyloidy. Two recent molecular phylogenetic studies have confirmed this classification as natural.
Description
A large, conspicuous mushroom, Amanita muscaria is generally common and numerous where it grows, and is often found in groups with basidiocarps in all stages of development.
Fly agaric fruiting bodies emerge from the soil looking like white eggs. After emerging from the ground, the cap is covered with numerous small white to yellow pyramid-shaped warts. These are remnants of the universal veil, a membrane that encloses the entire mushroom when it is still very young. Dissecting the mushroom at this stage reveals a characteristic yellowish layer of skin under the veil, which helps identification. As the fungus grows, the red colour appears through the broken veil and the warts become less prominent; they do not change in size, but are reduced relative to the expanding skin area. The cap changes from globose to hemispherical, and finally to plate-like and flat in mature specimens. Fully grown, the bright red cap is usually around in diameter, although larger specimens have been found. The red colour may fade after rain and in older mushrooms.
The free gills are white, as is the spore print. The oval spores measure 9–13 by 6.5–9 μm; they do not turn blue with the application of iodine. The stipe is white, high by wide, and has the slightly brittle, fibrous texture typical of many large mushrooms. At the base is a bulb that bears universal veil remnants in the form of two to four distinct rings or ruffs. Between the basal universal veil remnants and gills are remnants of the partial veil (which covers the gills during development) in the form of a white ring. It can be quite wide and flaccid with age. There is generally no associated smell other than a mild earthiness.
Although very distinctive in appearance, the fly agaric has been mistaken for other yellow to red mushroom species in the Americas, such as Armillaria cf. mellea and the edible A. basii—a Mexican species similar to A. caesarea of Europe. Poison control centres in the U.S. and Canada have become aware that (Spanish for 'yellow') is a common name for the A. caesarea-like species in Mexico. A. caesarea is distinguished by its entirely orange to red cap, which lacks the numerous white warty spots of the fly agaric (though these sometimes wash away during heavy rain). Furthermore, the stem, gills and ring of A. caesarea are bright yellow, not white. The volva is a distinct white bag, not broken into scales. In Australia, the introduced fly agaric may be confused with the native vermilion grisette (Amanita xanthocephala), which grows in association with eucalypts. The latter species generally lacks the white warts of A. muscaria and bears no ring. Additionally, immature button forms resemble puffballs.
Controversy
Amanita muscaria varies considerably in its morphology, and many authorities recognize several subspecies or varieties within the species. In The Agaricales in Modern Taxonomy, German mycologist Rolf Singer listed three subspecies, though without description: A. muscaria ssp. muscaria, A. muscaria ssp. americana, and A. muscaria ssp. flavivolvata.
However, a 2006 molecular phylogenetic study of different regional populations of A. muscaria by mycologist József Geml and colleagues found three distinct clades within this species representing, roughly, Eurasian, Eurasian "subalpine", and North American populations. Specimens belonging to all three clades have been found in Alaska; this has led to the hypothesis that this was the centre of diversification for this species. The study also looked at four named varieties of the species: var. alba, var. flavivolvata, var. formosa (including var. guessowii), and var. regalis from both areas. All four varieties were found within both the Eurasian and North American clades, evidence that these morphological forms are polymorphisms rather than distinct subspecies or varieties. Further molecular study by Geml and colleagues published in 2008 show that these three genetic groups, plus a fourth associated with oak–hickory–pine forest in the southeastern United States and two more on Santa Cruz Island in California, are delineated from each other enough genetically to be considered separate species. Thus A. muscaria as it stands currently is, evidently, a species complex. The complex also includes at least three other closely related taxa that are currently regarded as species: A. breckonii is a buff-capped mushroom associated with conifers from the Pacific Northwest, and the brown-capped A. gioiosa and A. heterochroma from the Mediterranean Basin and from Sardinia respectively. Both of these last two are found with Eucalyptus and Cistus trees, and it is unclear whether they are native or introduced from Australia.
Amanitaceae.org lists four varieties , but says that they will be segregated into their own taxa "in the near future". They are:
Distribution and habitat
A. muscaria is a cosmopolitan mushroom, native to conifer and deciduous woodlands throughout the temperate and boreal regions of the Northern Hemisphere, including higher elevations of warmer latitudes in regions such as Hindu Kush, the Mediterranean and also Central America. A recent molecular study proposes that it had an ancestral origin in the Siberian–Beringian region in the Tertiary period, before radiating outwards across Asia, Europe and North America. The season for fruiting varies in different climates: fruiting occurs in summer and autumn across most of North America, but later in autumn and early winter on the Pacific coast. This species is often found in similar locations to Boletus edulis, and may appear in fairy rings. Conveyed with pine seedlings, it has been widely transported into the southern hemisphere, including Australia, New Zealand, South Africa and South America, where it can be found in the Brazilian states of Paraná, São Paulo, Minas Gerais, Rio Grande do Sul.
Ectomycorrhizal, A. muscaria forms symbiotic relationships with many trees, including pine, oak, spruce, fir, birch, and cedar. Commonly seen under introduced trees, A. muscaria is the fungal equivalent of a weed in New Zealand, Tasmania and Victoria, forming new associations with southern beech (Nothofagus). The species is also invading a rainforest in Australia, where it may be displacing the native species. It appears to be spreading northwards, with recent reports placing it near Port Macquarie on the New South Wales north coast. It was recorded under silver birch (Betula pendula) in Manjimup, Western Australia in 2010. Although it has apparently not spread to eucalypts in Australia, it has been recorded associating with them in Portugal. Commonly found throughout the great Southern region of western Australia, it is regularly found growing on Pinus radiata.
Toxicity
A. muscaria poisoning has occurred in young children and in people who ingested the mushrooms for a hallucinogenic experience, or who confused it with an edible species.
A. muscaria contains several biologically active agents, at least one of which, muscimol, is known to be psychoactive. Ibotenic acid, a neurotoxin, serves as a prodrug to muscimol, with a small amount likely converting to muscimol after ingestion. An active dose in adults is approximately 6 mg muscimol or 30 to 60 mg ibotenic acid; this is typically about the amount found in one cap of Amanita muscaria. The amount and ratio of chemical compounds per mushroom varies widely from region to region and season to season, which can further confuse the issue. Spring and summer mushrooms have been reported to contain up to 10 times more ibotenic acid and muscimol than autumn fruitings.
Deaths from A. muscaria have been reported in historical journal articles and newspaper reports, but with modern medical treatment, fatal poisoning from ingesting this mushroom is extremely rare. Many books list A. muscaria as deadly, but according to David Arora, this is an error that implies the mushroom is far more toxic than it is. Furthermore, The North American Mycological Association has stated that there were "no reliably documented cases of death from toxins in these mushrooms in the past 100 years".
The active constituents of this species are water-soluble, and boiling and then discarding the cooking water at least partly detoxifies A. muscaria. Drying may increase potency, as the process facilitates the conversion of ibotenic acid to the more potent muscimol. According to some sources, once detoxified, the mushroom becomes edible. Patrick Harding describes the Sami custom of processing the fly agaric through reindeer.
Pharmacology
Muscarine, discovered in 1869, was long thought to be the active hallucinogenic agent in A. muscaria. Muscarine binds with muscarinic acetylcholine receptors leading to the excitation of neurons bearing these receptors. The levels of muscarine in Amanita muscaria are minute when compared with other poisonous fungi such as Inosperma erubescens, the small white Clitocybe species C. dealbata and C. rivulosa. The level of muscarine in A. muscaria is too low to play a role in the symptoms of poisoning.
The major toxins involved in A. muscaria poisoning are muscimol (3-hydroxy-5-aminomethyl-1-isoxazole, an unsaturated cyclic hydroxamic acid) and the related amino acid ibotenic acid. Muscimol is the product of the decarboxylation (usually by drying) of ibotenic acid. Muscimol and ibotenic acid were discovered in the mid-20th century. Researchers in England, Japan, and Switzerland showed that the effects produced were due mainly to ibotenic acid and muscimol, not muscarine. These toxins are not distributed uniformly in the mushroom. Most are detected in the cap of the fruit, a moderate amount in the base, with the smallest amount in the stalk. Quite rapidly, between 20 and 90 minutes after ingestion, a substantial fraction of ibotenic acid is excreted unmetabolised in the urine of the consumer. Almost no muscimol is excreted when pure ibotenic acid is eaten, but muscimol is detectable in the urine after eating A. muscaria, which contains both ibotenic acid and muscimol.
Ibotenic acid and muscimol are structurally related to each other and to two major neurotransmitters of the central nervous system: glutamic acid and GABA respectively. Ibotenic acid and muscimol act like these neurotransmitters, muscimol being a potent GABAA agonist, while ibotenic acid is an agonist of NMDA glutamate receptors and certain metabotropic glutamate receptors which are involved in the control of neuronal activity. It is these interactions which are thought to cause the psychoactive effects found in intoxication.
Muscazone is another compound that has more recently been isolated from European specimens of the fly agaric. It is a product of the breakdown of ibotenic acid by ultraviolet radiation. Muscazone is of minor pharmacological activity compared with the other agents. Amanita muscaria and related species are known as effective bioaccumulators of vanadium; some species concentrate vanadium to levels of up to 400 times those typically found in plants. Vanadium is present in fruit-bodies as an organometallic compound called amavadine. The biological importance of the accumulation process is unknown.
Symptoms
Fly agarics are best known for the unpredictability of their effects. Depending on habitat and the amount ingested per body weight, effects can range from mild nausea and twitching to drowsiness, cholinergic crisis-like effects (low blood pressure, sweating and salivation), auditory and visual distortions, mood changes, euphoria, relaxation, ataxia, and loss of equilibrium (like with tetanus.)
In cases of serious poisoning the mushroom causes delirium, somewhat similar in effect to anticholinergic poisoning (such as that caused by Datura stramonium), characterised by bouts of marked agitation with confusion, hallucinations, and irritability followed by periods of central nervous system depression. Seizures and coma may also occur in severe poisonings. Symptoms typically appear after around 30 to 90 minutes and peak within three hours, but certain effects can last for several days. In the majority of cases recovery is complete within 12 to 24 hours. The effect is highly variable between individuals, with similar doses potentially causing quite different reactions. Some people suffering intoxication have exhibited headaches up to ten hours afterwards. Retrograde amnesia and somnolence can result following recovery.
Treatment
Medical attention should be sought in cases of suspected poisoning. If the delay between ingestion and treatment is less than four hours, activated charcoal is given. Gastric lavage can be considered if the patient presents within one hour of ingestion. Inducing vomiting with syrup of ipecac is no longer recommended in any poisoning situation.
There is no antidote, and supportive care is the mainstay of further treatment for intoxication. Though sometimes referred to as a deliriant and while muscarine was first isolated from A. muscaria and as such is its namesake, muscimol does not have action, either as an agonist or antagonist, at the muscarinic acetylcholine receptor site, and therefore atropine or physostigmine as an antidote is not recommended. If a patient is delirious or agitated, this can usually be treated by reassurance and, if necessary, physical restraints. A benzodiazepine such as diazepam or lorazepam can be used to control combativeness, agitation, muscular overactivity, and seizures. Only small doses should be used, as they may worsen the respiratory depressant effects of muscimol. Recurrent vomiting is rare, but if present may lead to fluid and electrolyte imbalances; intravenous rehydration or electrolyte replacement may be required. Serious cases may develop loss of consciousness or coma, and may need intubation and artificial ventilation. Hemodialysis can remove the toxins, although this intervention is generally considered unnecessary. With modern medical treatment the prognosis is typically good following supportive treatment.
Uses
Psychoactive
The wide range of psychoactive effects have been variously described as depressant, sedative-hypnotic, psychedelic, dissociative, or deliriant; paradoxical effects such as stimulation may occur however. Perceptual phenomena such as synesthesia, macropsia, and micropsia may occur; the latter two effects may occur either simultaneously or alternatingly, as part of Alice in Wonderland syndrome, collectively known as dysmetropsia, along with related distortions pelopsia and teleopsia. Some users report lucid dreaming under the influence of its hypnotic effects. Unlike Psilocybe cubensis, A. muscaria cannot be commercially cultivated, due to its mycorrhizal relationship with the roots of pine trees. However, following the outlawing of psilocybin mushrooms in the United Kingdom in 2006, the sale of the still legal A. muscaria began increasing.
Marija Gimbutas reported to R. Gordon Wasson that in remote areas of Lithuania, A. muscaria has been consumed at wedding feasts, in which mushrooms were mixed with vodka. She also reported that the Lithuanians used to export A. muscaria to the Sami in the Far North for use in shamanic rituals. The Lithuanian festivities are the only report that Wasson received of ingestion of fly agaric for religious use in Eastern Europe.
Siberia
A. muscaria was widely used as an entheogen by many of the indigenous peoples of Siberia. Its use was known among almost all of the Uralic-speaking peoples of western Siberia and the Paleosiberian-speaking peoples of the Russian Far East. There are only isolated reports of A. muscaria use among the Tungusic and Turkic peoples of central Siberia and it is believed that on the whole entheogenic use of A. muscaria was not practised by these peoples. In western Siberia, the use of A. muscaria was restricted to shamans, who used it as an alternative method of achieving a trance state. (Normally, Siberian shamans achieve trance by prolonged drumming and dancing.) In eastern Siberia, A. muscaria was used by both shamans and laypeople alike, and was used recreationally as well as religiously. In eastern Siberia, the shaman would take the mushrooms, and others would drink his urine. This urine, still containing psychoactive elements, may be more potent than the A. muscaria mushrooms with fewer negative effects such as sweating and twitching, suggesting that the initial user may act as a screening filter for other components in the mushroom.
The Koryak of eastern Siberia have a story about the fly agaric (wapaq) which enabled Big Raven to carry a whale to its home. In the story, the deity Vahiyinin ("Existence") spat onto earth, and his spittle became the wapaq, and his saliva becomes the warts. After experiencing the power of the wapaq, Raven was so exhilarated that he told it to grow forever on earth so his children, the people, could learn from it. Among the Koryaks, one report said that the poor would consume the urine of the wealthy, who could afford to buy the mushrooms. It was reported that the local reindeer would often follow an individual intoxicated by the muscimol mushroom, and if said individual were to urinate in snow the reindeer would become similarly intoxicated and the Koryak people's would use the drunken state of the reindeer to more easily rope and hunt them.
Recent rise in popularity
As a result of a lack of regulation, the use of Amanita muscaria as a popular legal alternative to hallucinogens has grown exponentially in recent years. In 2024, Google searches for Amanita muscaria rose nearly 200% from the previous year, a trend that an article published in the American Journal of Preventative Medicine correlated with the sudden commercialization of Amanita muscaria products on the internet.
While Amanita mushrooms are unscheduled in the United States, the sale of Amanita products exists in a legal gray area as they are listed as a poison by the FDA and are not approved to be used in dietary supplements, with some drawing comparisons to the controversial legal status of hemp-derived cannabinoids.
A recent outbreak of poisonings and at least one death associated with products containing Amanita muscaria extracts has sparked debates regarding the regulatory status of Amanita mushrooms and their psychoactive constituents. These products often use misleading advertising, such as erroneous comparisons to Psilocybin mushrooms or simply not disclosing the inclusion of Amanita mushrooms on the packaging.
Other reports and theories
The Finnish historian T. I. Itkonen mentions that A. muscaria was once used among the Sámi peoples. Sorcerers in Inari would consume fly agarics with seven spots. In 1979, Said Gholam Mochtar and Hartmut Geerken published an article in which they claimed to have discovered a tradition of medicinal and recreational use of this mushroom among a Parachi-speaking group in Afghanistan. There are also unconfirmed reports of religious use of A. muscaria among two Subarctic Native American tribes. Ojibwa ethnobotanist Keewaydinoquay Peschel reported its use among her people, where it was known as (an abbreviation of the name (= "red-top mushroom"). This information was enthusiastically received by Wasson, although evidence from other sources was lacking. There is also one account of a Euro-American who claims to have been initiated into traditional Tlicho use of Amanita muscaria.
Mycophilosopher Martijn Benders has proposed a novel evolutionary theory involving Amanita muscaria. In his book Amanita Muscaria – the Book of the Empress, Benders argues that a precursor of ibotenic acid, a compound found in the mushroom, was present in ancient seaweed and played a significant role in the evolution of life. According to this hypothesis, the compound influenced the twitching movements of early aquatic organisms, leading to the development of behaviors such as jumping onto land—a crucial step in the evolution of terrestrial species.
The flying reindeer of Santa Claus, who is called Joulupukki in Finland, could symbolize the use of A. muscaria by Sámi shamans. However, Sámi scholars and the Sámi peoples themselves refute any connection between Santa Claus and Sámi history or culture."The story of Santa emerging from a Sámi shamanic tradition has a critical number of flaws," asserts Tim Frandy, assistant professor of Nordic Studies at the University of British Columbia and a member of the Sámi descendent community in North America. "The theory has been widely criticized by Sámi people as a stereotypical and problematic romanticized misreading of actual Sámi culture."
Vikings
The notion that Vikings used A. muscaria to produce their berserker rages was first suggested by the Swedish professor Samuel Ödmann in 1784. Ödmann based his theories on reports about the use of fly agaric among Siberian shamans. The notion has become widespread since the 19th century, but no contemporary sources mention this use or anything similar in their description of berserkers. Muscimol is generally a mild relaxant, but it can create a range of different reactions within a group of people. It is possible that it could make a person angry, or cause them to be "very jolly or sad, jump about, dance, sing or give way to great fright". Comparative analysis of symptoms have, however, since shown Hyoscyamus niger to be a better fit to the state that characterises the berserker rage.
Soma
In 1968, R. Gordon Wasson proposed that A. muscaria was the soma talked about in the Rigveda of India, a claim which received widespread publicity and popular support at the time. He noted that descriptions of Soma omitted any description of roots, stems or seeds, which suggested a mushroom, and used the adjective hári "dazzling" or "flaming" which the author interprets as meaning red. One line described men urinating Soma; this recalled the practice of recycling urine in Siberia. Soma is mentioned as coming "from the mountains", which Wasson interpreted as the mushroom having been brought in with the Aryan migrants from the north. Indian scholars Santosh Kumar Dash and Sachinanda Padhy pointed out that both eating of mushrooms and drinking of urine were proscribed, using as a source the Manusmṛti.
In 1971, Vedic scholar John Brough from Cambridge University rejected Wasson's theory and noted that the language was too vague to determine a description of Soma. In his 1976 survey, Hallucinogens and Culture, anthropologist Peter T. Furst evaluated the evidence for and against the identification of the fly agaric mushroom as the Vedic Soma, concluding cautiously in its favour. Kevin Feeney and Trent Austin compared the references in the Vedas with the filtering mechanisms in the preparation of Amanita muscaria and published findings supporting the proposal that fly-agaric mushrooms could be a likely candidate for the sacrament. Other proposed candidates include Psilocybe cubensis, Peganum harmala, and Ephedra.
Christianity
Philologist, archaeologist, and Dead Sea Scrolls scholar John Marco Allegro postulated that early Christian theology was derived from a fertility cult revolving around the entheogenic consumption of A. muscaria in his 1970 book The Sacred Mushroom and the Cross. This theory has found little support by scholars outside the field of ethnomycology. The book was widely criticized by academics and theologians, including Sir Godfrey Driver, emeritus Professor of Semitic Philology at Oxford University and Henry Chadwick, the Dean of Christ Church, Oxford. Christian author John C. King wrote a detailed rebuttal of Allegro's theory in the 1970 book A Christian View of the Mushroom Myth; he notes that neither fly agarics nor their host trees are found in the Middle East, even though cedars and pines are found there, and highlights the tenuous nature of the links between biblical and Sumerian names coined by Allegro. He concludes that if the theory were true, the use of the mushroom must have been "the best kept secret in the world" as it was so well concealed for two thousand years.
Fly trap
Amanita muscaria is traditionally used for catching flies possibly due to its content of ibotenic acid and muscimol, which lead to its common name "fly agaric". Recently, an analysis of nine different methods for preparing A. muscaria for catching flies in Slovenia have shown that the release of ibotenic acid and muscimol did not depend on the solvent (milk or water) and that thermal and mechanical processing led to faster extraction of ibotenic acid and muscimol.
Culinary
The toxins in A. muscaria are water-soluble: parboiling A. muscaria fruit bodies can detoxify them and render them edible, although consumption of the mushroom as a food has never been widespread. The consumption of detoxified A. muscaria has been practiced in some parts of Europe (notably by Russian settlers in Siberia) since at least the 19th century, and likely earlier. The German physician and naturalist Georg Heinrich von Langsdorff wrote the earliest published account on how to detoxify this mushroom in 1823. In the late 19th century, the French physician Félix Archimède Pouchet was a populariser and advocate of A. muscaria consumption, comparing it to manioc, an important food source in tropical South America that must also be detoxified before consumption.
Use of this mushroom as a food source also seems to have existed in North America. A classic description of this use of A. muscaria by an African-American mushroom seller in Washington, D.C., in the late 19th century is described by American botanist Frederick Vernon Coville. In this case, the mushroom, after parboiling, and soaking in vinegar, is made into a mushroom sauce for steak. It is also consumed as a food in parts of Japan. The most well-known current use as an edible mushroom is in Nagano Prefecture, Japan. There, it is primarily salted and pickled.
A 2008 paper by food historian William Rubel and mycologist David Arora gives a history of consumption of A. muscaria as a food and describes detoxification methods. They advocate that Amanita muscaria be described in field guides as an edible mushroom, though accompanied by a description on how to detoxify it. The authors state that the widespread descriptions in field guides of this mushroom as poisonous is a reflection of cultural bias, as several other popular edible species, notably morels, are also toxic unless properly cooked.
In culture
The red-and-white spotted toadstool is a common image in many aspects of popular culture. Garden ornaments and children's picture books depicting gnomes and fairies, such as the Smurfs, often show fly agarics used as seats, or homes. Fly agarics have been featured in paintings since the Renaissance, albeit in a subtle manner. For instance, in Hieronymus Bosch's painting, The Garden of Earthly Delights, the mushroom can be seen on the left-hand panel of the work. In the Victorian era they became more visible, becoming the main topic of some fairy paintings. Two of the most famous uses of the mushroom are in the Mario franchise (specifically two of the Super Mushroom power-up items and the platforms in several stages which are based on a fly agaric), and the dancing mushroom sequence in the 1940 Disney film Fantasia.
An account of the journeys of Philip von Strahlenberg to Siberia and his descriptions of the use of the mukhomor there was published in English in 1736. The drinking of urine of those who had consumed the mushroom was commented on by Anglo-Irish writer Oliver Goldsmith in his widely read 1762 novel, Citizen of the World. The mushroom had been identified as the fly agaric by this time. Other authors recorded the distortions of the size of perceived objects while intoxicated by the fungus, including naturalist Mordecai Cubitt Cooke in his books The Seven Sisters of Sleep and A Plain and Easy Account of British Fungi. This observation is thought to have formed the basis of the effects of eating the mushroom in the 1865 popular story Alice's Adventures in Wonderland. A hallucinogenic "scarlet toadstool" from Lappland is featured as a plot element in Charles Kingsley's 1866 novel Hereward the Wake based on the medieval figure of the same name. Thomas Pynchon's 1973 novel Gravity's Rainbow describes the fungus as a "relative of the poisonous Destroying angel" and presents a detailed description of a character preparing a cookie bake mixture from harvested Amanita muscaria. Fly agaric shamanism—in the context of a surviving Dionysian cult in the Peak District—is also explored in the 2003 novel Thursbitch by Alan Garner.
| Biology and health sciences | Poisonous fungi | null |
59399 | https://en.wikipedia.org/wiki/Lumber | Lumber | Lumber is wood that has been processed into uniform and useful sizes (dimensional lumber), including beams and planks or boards. Lumber is mainly used for construction framing, as well as finishing (floors, wall panels, window frames). Lumber has many uses beyond home building. Lumber is referred to as timber in the United Kingdom, Europe, Australia, and New Zealand, while in other parts of the world (mainly the United States and Canada) the term timber refers specifically to unprocessed wood fiber, such as cut logs or standing trees that have yet to be cut.
Lumber may be supplied either rough-sawn, or surfaced on one or more of its faces. Rough lumber is the raw material for furniture-making, and manufacture of other items requiring cutting and shaping. It is available in many species, including hardwoods and softwoods, such as white pine and red pine, because of their low cost.
Finished lumber is supplied in standard sizes, mostly for the construction industry – primarily softwood, from coniferous species, including pine, fir and spruce (collectively spruce-pine-fir), cedar, and hemlock, but also some hardwood, for high-grade flooring. It is more commonly made from softwood than hardwoods, and 80% of lumber comes from softwood.
Terminology
In the United States and Canada, milled boards are called lumber, while timber describes standing or felled trees.
In contrast, in Britain, and some other Commonwealth nations and Ireland, the term timber is used in both senses. (In the UK, the word lumber is rarely used in relation to wood and has several other meanings.)
Re-manufactured lumber
Re-manufactured lumber is the result of secondary or tertiary processing of previously milled lumber. Specifically, it refers to lumber cut for industrial or wood-packaging use. Lumber is cut by ripsaw or resaw to create dimensions that are not usually processed by a primary sawmill.
Re-sawing is the splitting of hardwood or softwood lumber into two or more thinner pieces of full-length boards. For example, splitting a 2×4 () into two 1×4s () of the same length is considered re-sawing.
Plastic lumber
Structural lumber may also be produced from recycled plastic and new plastic stock. Its introduction has been strongly opposed by the forestry industry. Blending fiberglass in plastic lumber enhances its strength, durability, and fire resistance. Plastic fiberglass structural lumber can have a "class 1 flame spread rating of 25 or less, when tested in accordance with ASTM standard E 84," which means it burns more slowly than almost all treated wood lumber.
Timber mark
A timber mark is a code beaten on to cut wood by a specially made hammer to show the logging licence.
History
The definition of the word lumber as sawn planks of wood originated in the 17th century in North America.
In 1420, the archipelago of Madeira was colonized by the Portuguese Empire. Prince Henry the Navigator sent settlers to Madeira, who cleared the huge expanses of forest to grow crops. The felled trees were processed at sawmills and shipped to the mainland.
Cornelis Corneliszoon (or Krelis Lootjes) was a Dutch windmill owner from Uitgeest who invented the first wind-powered sawmill in 1593. This made the conversion of logs into planks thirty times faster than previous manually operated sawmills.
Conversion of wood logs
Logs are converted into lumber by being sawn, hewn, or split. Sawing with a rip saw is the most common method, because sawing allows logs of lower quality, with irregular grain and large knots, to be used and is more economical. There are various types of sawing:
Plain sawn (flat sawn, through and through, bastard sawn) – A log sawn through without adjusting the position of the log and the grain runs across the width of the boards.
Quarter sawn and rift sawn – These terms have been confused in history but generally mean lumber sawn so the annual rings are reasonably perpendicular to the sides (not edges) of the lumber.
Boxed heart – The pith remains within the timber, post or beam, with some allowance for exposure.
Heart center – the center core of a log.
Free of heart center (FOHC) – A side-cut timber, post or beam without any pith.
Free of knots (FOK) – No knots are present.
Dimensional lumber
Dimensional lumber is lumber that is cut to standardized width and depth, often specified in millimetres or inches (but see below for information on nominal dimensions vs. actual dimensions). Carpenters extensively use dimensional lumber in framing wooden buildings. Common sizes include 2×4 (pictured) (also two-by-four and other variants, such as four-by-two in Australia, New Zealand, and the UK), 2×6, and 4×4. The length of a board is usually specified separately from the width and depth. It is thus possible to find 2×4s that are four, eight, and twelve feet in length. In Canada and the United States, the standard lengths of lumber are . For wall framing, precut "stud" lengths are available, and are commonly used. For ceilings heights of , studs are available in , , and .
North American softwoods
The length of a unit of dimensional lumber is limited by the height and girth of the tree it is milled from. In general the maximum length is . Engineered wood products, manufactured by binding the strands, particles, fibers, or veneers of wood, together with adhesives, to form composite materials, offer more flexibility and greater structural strength than typical wood building materials.
Pre-cut studs save a framer much time, because they are pre-cut by the manufacturer for use in 8-, 9-, and 10-foot ceiling applications, which means the manufacturer has removed a few inches or centimetres of the piece to allow for the sill plate and the double top plate with no additional sizing necessary.
In the Americas, two-bys (2×4s, 2×6s, 2×8s, 2×10s, and 2×12s), named for traditional board thickness in inches, along with the 4×4 (), are common lumber sizes used in modern construction. They are the basic building blocks for such common structures as balloon-frame or platform-frame housing. Dimensional lumber made from softwood is typically used for construction, while hardwood boards are more commonly used for making cabinets or furniture.
Lumber's nominal dimensions are larger than the actual standard dimensions of finished lumber. Historically, the nominal dimensions were the size of the green (not dried), rough (unfinished) boards that eventually became smaller finished lumber through drying and planing (to smooth the wood). Today, the standards specify the final finished dimensions and the mill cuts the logs to whatever size it needs to achieve those final dimensions. Typically, that rough cut is smaller than the nominal dimensions because modern technology makes it possible to use the logs more efficiently. For example, a "2×4" board historically started out as a green, rough board actually . After drying and planing, it would be smaller by a nonstandard amount. Today, a "2×4" board starts out as something smaller than 2 inches by 4 inches and not specified by standards, and after drying and planing is minimally .
As previously noted, less wood is needed to produce a given finished size than when standards called for the green lumber to be the full nominal dimension. However, even the dimensions for finished lumber of a given nominal size have changed over time. In 1910, a typical finished board was . In 1928, that was reduced by 4%, and yet again by 4% in 1956. In 1961, at a meeting in Scottsdale, Arizona, the Committee on Grade Simplification and Standardization agreed to what is now the current U.S. standard: in part, the dressed size of a 1-inch (nominal) board was fixed at inch; while the dressed size of 2 inch (nominal) lumber was reduced from inch to the current inch.
In 1964, Popular Mechanics magazine hired an independent agency to test the comparative strength of multiple samples of (A) a full-size 2×4 inches, (B) × inches, (C) × inches, and (D) × inches (today’s standard). With A’s compressive strength benchmarked as “100%,” B-C-D were 90.7%, 82.2%, and 73.6% the strength of A’s full-size 2×4. Stated another way, the 1960s’ reduction of the smaller dimension from to inches reduced compressive strength by 10.46%.
Dimensional lumber is available in green, unfinished state, and for that kind of lumber, the nominal dimensions are the actual dimensions.
Grades and standards
Individual pieces of lumber exhibit a wide range in quality and appearance with respect to knots, slope of grain, shakes and other natural characteristics. Therefore, they vary considerably in strength, utility, and value.
The move to set national standards for lumber in the United States began with the publication of the American Lumber Standard in 1924, which set specifications for lumber dimensions, grade, and moisture content; it also developed inspection and accreditation programs. These standards have changed over the years to meet the changing needs of manufacturers and distributors, with the goal of keeping lumber competitive with other construction products. Current standards are set by the American Lumber Standard Committee, appointed by the U.S. Secretary of Commerce.
Design values for most species and grades of visually graded structural products are determined in accordance with ASTM standards, which consider the effect of strength reducing characteristics, load duration, safety, and other influencing factors. The applicable standards are based on results of tests conducted in cooperation with the USDA Forest Products Laboratory. Design Values for Wood Construction, which is a supplement to the ANSI/AF&PA National Design Specification® for Wood Construction, provides these lumber design values, which are recognized by the model building codes.
Canada has grading rules that maintain a standard among mills manufacturing similar woods to assure customers of uniform quality. Grades standardize the quality of lumber at different levels and are based on moisture content, size, and manufacture at the time of grading, shipping, and unloading by the buyer. The National Lumber Grades Authority (NLGA) is responsible for writing, interpreting and maintaining Canadian lumber grading rules and standards. The Canadian Lumber Standards Accreditation Board (CLSAB) monitors the quality of Canada's lumber grading and identification system. Their common grade abbrievation, CLS, Canadian Lumber Standard is well utilised in the construction industry.
Attempts to maintain lumber quality over time have been challenged by historical changes in the timber resources of the United States – from the slow-growing virgin forests common over a century ago to the fast-growing plantations now common in today's commercial forests. Resulting declines in lumber quality have been of concern to both the lumber industry and consumers and have caused increased use of alternative construction products.
Machine stress-rated and machine-evaluated lumber are readily available for end-uses where high strength is critical, such as trusses, rafters, laminating stock, I-beams and web joints. Machine grading measures a characteristic such as stiffness or density that correlates with the structural properties of interest, such as bending strength. The result is a more precise understanding of the strength of each piece of lumber than is possible with visually graded lumber, which allows designers to use full-design strength and avoid overbuilding.
In Europe, strength grading of rectangular sawn lumber/timber (both softwood and hardwood) is done according to EN-14081 and commonly sorted into classes defined by EN-338. For softwoods, the common classes are (in increasing strength) C16, C18, C24, and C30. There are also classes specifically for hardwoods and those in most common use (in increasing strength) are D24, D30, D40, D50, D60, and D70. For these classes, the number refers to the required 5th percentile bending strength in newtons per square millimetre. There are other strength classes, including T-classes based on tension intended for use in glulam.
C14, used for scaffolding and formwork
C16 and C24, general construction
C30, prefab roof trusses and where design requires somewhat stronger joists than C24 can offer. TR26 is also a common trussed rafter strength class in long standing use in the UK.
C40, usually seen in glulam
Grading rules for African and South American sawn lumber have been developed by ATIBT according to the rules of the Sciages Avivés Tropicaux Africains (SATA) and is based on clear cuttings – established by the percentage of the clear surface.
North American hardwoods
In North America, market practices for dimensional lumber made from hardwoods varies significantly from the regularized standardized 'dimension lumber' sizes used for sales and specification of softwoods – hardwood boards are often sold totally rough cut, or machine planed only on the two (broader) face sides. When hardwood boards are also supplied with planed faces, it is usually both by random widths of a specified thickness (normally matching milling of softwood dimensional lumber) and somewhat random lengths. But besides those older (traditional and normal) situations, in recent years some product lines have been widened to also market boards in standard stock sizes; these usually retail in big-box stores and using only a relatively small set of specified lengths; in all cases hardwoods are sold to the consumer by the board-foot (), whereas that measure is not used for softwoods at the retailer (to the cognizance of the buyer).
Also in North America, hardwood lumber is commonly sold in a "quarter" system, when referring to thickness; 4/4 (four quarter) refers to a board, 8/4 (eight quarter) is a board, etc. This "quarter" system is rarely used for softwood lumber; although softwood decking is sometimes sold as 5/4, even though it is actually one inch thick (from milling off each side in a motorized planing step of production). The "quarter" system of reference is a traditional North American lumber industry nomenclature used specifically to indicate the thickness of rough sawn hardwood lumber.
In rough-sawn lumber it immediately clarifies that the lumber is not yet milled, avoiding confusion with milled dimension lumber which is measured as actual thickness after machining. Examples – -inch, 19 mm, or 1x. In recent years architects, designers, and builders have begun to use the "quarter" system in specifications as a vogue of insider knowledge, though the materials being specified are finished lumber, thus conflating the separate systems and causing confusion.
Hardwoods cut for furniture are cut in the fall and winter, after the sap has stopped running in the trees. If hardwoods are cut in the spring or summer the sap ruins the natural color of the lumber and decreases the value of the wood for furniture.
Engineered lumber
Engineered lumber is lumber created by a manufacturer and designed for a certain structural purpose. The main categories of engineered lumber are:
Laminated veneer lumber (LVL) – LVL comes in thicknesses with depths such as , and are often doubled or tripled up. They function as beams to provide support over large spans, such as removed support walls and garage door openings, places where dimensional lumber is insufficient, and also in areas where a heavy load is bearing from a floor, wall or roof above on a somewhat short span where dimensional lumber is impractical. This type of lumber is compromised if it is altered by holes or notches anywhere within the span or at the ends, but nails can be driven into it wherever necessary to anchor the beam or to add hangers for I-joists or dimensional lumber joists that terminate at an LVL beam.
Wooden I-joists – sometimes called "TJI", "Trus Joists" or "BCI", all of which are brands of wooden I-joists, they are used for floor joists on upper floors and also in first floor conventional foundation construction on piers as opposed to slab floor construction. They are engineered for long spans and are doubled up in places where a wall will be aligned over them, and sometimes tripled where heavy roof-loaded support walls are placed above them. They consist of a top and bottom chord or flange made from dimensional lumber with a webbing in-between made from oriented strand board (OSB) (or, latterly, steel mesh forms which allow passage of services without cutting). The webbing can be removed up to certain sizes or shapes according to the manufacturer's or engineer's specifications, but for small holes, wooden I-joists come with "knockouts", which are perforated, pre-cut areas where holes can be made easily, typically without engineering approval. When large holes are needed, they can typically be made in the webbing only and only in the center third of the span; the top and bottom chords lose their integrity if cut. Sizes and shapes of the hole, and typically the placing of a hole itself, must be approved by an engineer prior to the cutting of the hole and in many areas, a sheet showing the calculations made by the engineer must be provided to the building inspection authorities before the hole will be approved. Some I-joists are made with W-style webbing like a truss to eliminate cutting and to allow ductwork to pass through.
Finger-jointed lumber – solid dimensional lumber lengths typically are limited to lengths of , but can be made longer by the technique of "finger-jointing" by using small solid pieces, usually long, and joining them together using finger joints and glue to produce lengths that can be up to long in 2×6 size. Finger-jointing also is predominant in precut wall studs. It is also an affordable alternative for non-structural hardwood that will be painted (staining would leave the finger-joints visible). Care is taken during construction to avoid nailing directly into a glued joint as stud breakage can occur.
Glulam beams – created from 2×4 or 2×6 stock by gluing the faces together to create beams such as 4×12 or 6×16. As such, a beam acts as one larger piece of lumber – thus eliminating the need to harvest larger, older trees for the same size beam.
Manufactured trusses – trusses are used in home construction as a pre-fabricated replacement for roof rafters and ceiling joists (stick-framing). It is seen as an easier installation and a better solution for supporting roofs than the use of dimensional lumber's struts and purlins as bracing. In the southern U.S. and elsewhere, stick-framing with dimensional lumber roof support is still predominant. The main drawbacks of trusses are reduced attic space, time required for engineering and ordering, and a cost higher than the dimensional lumber needed if the same project were conventionally framed. The advantages are significantly reduced labor costs (installation is faster than conventional framing), consistency, and overall schedule savings.
Various pieces and cuts
Square and rectangular forms: plank, slat, batten, board, lath, strapping (typically ), cant (A partially sawn log such as sawn on two sides or squared to a large size and later resawn into lumber. A flitch is a type of cant with wane on one or both sides). Various pieces are also known by their uses such as post, beam, (girt), stud, rafter, joist, sill plate, wall plate.
Rod forms: pole, (dowel), stick (staff, baton)
Timber piles
In the United States, pilings are mainly cut from southern yellow pines and Douglas-fir. Treated pilings are available in chromated copper arsenate retentions of if treatment is required.
Historical Chinese construction
Under the prescription of the Method of Construction (營造法式) issued by the Song dynasty government in the early twelfth century, timbers were standardized to eight cross-sectional dimensions. Regardless of the actual dimensions of the timber, the ratio between width and height was maintained at 1:1.5. Units are in Song dynasty inches (31.2 mm).
Timber smaller than the 8th class were called "unclassed" (等外). The width of a timber is referred to as one "timber" (材), and the dimensions of other structural components were quoted in multiples of "timber"; thus, as the width of the actual timber varied, the dimensions of other components were easily calculated, without resorting to specific figures for each scale. The dimensions of timbers in similar applications show a gradual diminution from the Sui dynasty (580–618) to the modern era; a 1st class timber during the Sui was reconstructed as 15×10 (Sui dynasty inches, or 29.4 mm).
Defects in lumber
Defects occurring in lumber are grouped into the following four divisions:
Conversion
During the process of converting timber to commercial forms of lumber the following defects may occur:
Chip mark: this defect is indicated by the marks or signs placed by chips on the finished surface of timber
Diagonal grain: improper sawing of timber
Torn grain: when a small dent is made on the finished surface due to falling of some tool
Wane: presence of original rounded surface in the finished product
Defects due to fungi and animals
Fungi attack wood (both timber and lumber) when these conditions are all present:
The wood moisture content is above 25% on a dry-weight basis
The environment is sufficiently warm
Oxygen (O2) is present
Wood with less than 25% moisture (dry weight basis) can remain free of decay for centuries. Similarly, wood submerged in water may not be attacked by fungi if the amount of oxygen is inadequate.
Fungi lumber/timber defects:
Blue stain
Brown rot
Dry rot
Heart rot
Sap stain
Wet rot
White rot
Following are the insects and molluscs which are usually responsible for the decay of timber/lumber:
Woodboring beetles
Marine borers (Barnea similis)
Teredos (Teredo navalis)
Termites
Carpenter ants
Carpenter bees
Natural forces
There are two main natural forces responsible for causing defects in timber and lumber: abnormal growth and rupture of tissues. Rupture of tissue includes cracks or splits in the wood called "shakes". "Ring shake", "wind shake", or "ring failure" is when the wood grain separates around the growth rings either while standing or during felling. Shakes may reduce the strength of a timber and the appearance thus reduce lumber grade and may capture moisture, promoting decay. Eastern hemlock is known for having ring shake. A "check" is a crack on the surface of the wood caused by the outside of a timber shrinking as it seasons. Checks may extend to the pith and follow the grain. Like shakes, checks can hold water promoting rot. A "split" goes all the way through a timber. Checks and splits occur more frequently at the ends of lumber because of the more rapid drying in these locations.
Next to defects, uneven expansion or contraction caused by changes in moisture content will cause sawed timber to warp, making it less suitable for many purposes.
Seasoning
The seasoning of lumber is typically either kiln- or air-dried. Defects due to seasoning are the main cause of splits, bowing and honeycombing. Seasoning is the process of drying timber to remove the bound moisture contained in the walls of the wood cells to produce seasoned timber.
Durability and service life
Under proper conditions, wood provides excellent, lasting performance. However, it also faces several potential threats to service life, including fungal activity and insect damage – which can be avoided in numerous ways. Section 2304.11 of the International Building Code addresses protection against decay and termites. This section provides requirements for non-residential construction applications, such as wood used above ground (e.g., for framing, decks, stairs, etc.), as well as other applications.
There are four recommended methods to protect wood-frame structures against durability hazards and thus provide maximum service life for the building. All require proper design and construction:
Controlling moisture using design techniques to avoid decay
Providing effective control of termites and other insects
Using durable materials such as pressure-treated or naturally durable species of wood where appropriate
Providing quality assurance during design and construction and throughout the building's service life using appropriate maintenance practices
Moisture control
Wood is a hygroscopic material, which means it naturally absorbs and releases water to balance its internal moisture content with the surrounding environment. The moisture content of wood is measured by the weight of water as a percentage of the oven-dry weight of the wood fiber. The key to controlling decay is controlling moisture. Once decay fungi are established, the minimum moisture content for decay to propagate is 22 to 24 percent, so building experts recommend 19 percent as the maximum safe moisture content for untreated wood in service. Water by itself does not harm the wood, but rather, wood with consistently high moisture content enables fungal organisms to grow.
The primary objective when addressing moisture loads is to keep water from entering the building envelope in the first place and to balance the moisture content within the building itself. Moisture control by means of accepted design and construction details is a simple and practical method of protecting a wood-frame building against decay. For applications with a high risk of staying wet, designers specify durable materials such as naturally decay-resistant species or wood that has been treated with preservatives. Cladding, shingles, sill plates and exposed timbers or glulam beams are examples of potential applications for treated wood.
Controlling termites and other insects
For buildings in termite zones, basic protection practices addressed in current building codes include (but are not limited to) the following:
Grading the building site away from the foundation to provide proper drainage
Covering exposed ground in any crawl spaces with 6-mil polyethylene film and maintaining at least of clearance between the ground and the bottom of framing members above (12 inches to beams or girders, 18 inches to joists or plank flooring members)
Supporting post columns by concrete piers so that there is at least of clear space between the wood and exposed earth
Installing wood framing and sheathing in exterior walls at least eight inches above exposed earth; locating siding at least six inches from the finished grade
Where appropriate, ventilating crawl spaces according to local building codes
Removing building material scraps from the job site before backfilling.
If allowed by local regulation, treating the soil around the foundation with an approved termiticide to provide protection against subterranean termites
Preservatives
To avoid decay and termite infestation, untreated wood is separated from the ground and other sources of moisture. These separations are required by many building codes and are considered necessary to maintain wood elements in permanent structures at a safe moisture content for decay protection. When it is not possible to separate wood from the sources of moisture, designers often rely on preservative-treated wood.
Wood can be treated with a preservative that improves service life under severe conditions without altering its basic characteristics. It can also be pressure-impregnated with fire-retardant chemicals that improve its performance in a fire. One of the early treatments to "fireproof lumber", which retard fires, was developed in 1936 by the Protexol Corporation, in which lumber is heavily treated with salt. Wood does not deteriorate simply because it gets wet. When wood breaks down, it is because an organism is eating it. Preservatives work by making the food source inedible to these organisms. Properly preservative-treated wood can have 5 to 10 times the service life of untreated wood. Preserved wood is used most often for railroad ties, utility poles, marine piles, decks, fences and other outdoor applications. Various treatment methods and types of chemicals are available, depending on the attributes required in the particular application and the level of protection needed.
There are two basic methods of treating: with and without pressure. Non-pressure methods are the application of preservatives by brushing, spraying, or dipping the piece to be treated. Deeper, more thorough penetration is achieved by driving the preservative into the wood cells with pressure. Various combinations of pressure and vacuum are used to force adequate levels of chemical into the wood. Pressure-treating preservatives consist of chemicals carried in a solvent.
Chromated copper arsenate, once the most commonly used wood preservative in North America began being phased out of most residential applications in 2004. Replacing it are amine copper quat and copper azole.
All wood preservatives used in the United States and Canada are registered and regularly re-examined for safety by the U.S. Environmental Protection Agency and Health Canada's Pest Management and Regulatory Agency, respectively.
Timber framing
Timber framing is a style of construction that uses heavier framing elements (larger posts and beams) than modern stick framing, which uses smaller standard dimensional lumber. The timbers are cut from log boles and squared with a saw, broadaxe or adze, and then joined together with joinery without nails. Modern timber framing has been growing in popularity in the United States since the 1970s.
Environmental effects of lumber
Green building minimizes the impact or "environmental footprint" of a building. Wood is a major building material that is renewable and replenishable in a continuous cycle. Studies show manufacturing wood uses less energy and results in less air and water pollution than steel and concrete. However, demand for lumber is blamed for deforestation.
Residual wood
The conversion from coal to biomass power is a growing trend in the United States.
The United Kingdom, Uzbekistan, Kazakhstan, Australia, Fiji, Madagascar, Mongolia, Russia, Denmark, Switzerland, and Eswatini governments all support an increased role for energy derived from biomass, which are organic materials available on a renewable basis and include residues and/or byproducts of the logging, saw milling and paper-making processes. In particular, they view it as a way to lower greenhouse gas emissions by reducing the consumption of oil and gas while supporting the growth of forestry, agriculture and rural economies. Studies by the U.S. government have found the country's combined forest and agriculture land resources have the power to sustainably supply more than one-third of its current petroleum consumption.
Biomass is already an important source of energy for the North American forest products industry. It is common for companies to have cogeneration facilities, also known as combined heat and power, which convert some of the biomass that results from wood and paper manufacturing to electrical and thermal energy in the form of steam. The electricity is used to, among other things, dry lumber and supply heat to the dryers used in paper-making.
Environmental impacts
Lumber is a sustainable and environmentally friendly construction material that could replace modern building materials (e.g. concrete and steel) given its structural performance, capacity to fixate CO2 and low energy demand during the manufacturing process.
Substituting lumber for concrete or steel avoids the carbon emissions of those materials. Cement and concrete manufacture is responsible for around 8% of global GHG emissions while the iron and steel industry is responsible for another 5% (half a ton of CO2 is emitted to manufacture a ton of concrete; two tons of CO2 are emitted in the manufacture of a ton of steel).
Advantages of lumber:
Fire performance: In the case of fire, the outer layer of mass timber will tend to char in a predictable way that effectively self-extinguishes and shields the interior, allowing it to retain structural integrity for several hours, even in an intense fire.
Reduction of carbon emissions: Building materials and construction make up 11% of global greenhouse gas emissions. Though the exact amount will depend on tree species, forestry practices, transportation costs, and several other factors, that one cubic meter of lumber sequesters roughly one tonne of CO2. It is estimated that wood can reduce the amount of CO2 released into the atmosphere by half. In addition, wood has a significant CO2 storage capacity, which limits its release. However, when wood is destroyed (naturally or by combustion), all of the previously stored CO2 is released into the atmosphere.
Natural insulation: lumber is a natural insulator which makes it particularly good for windows and doors.
Less construction time, labor costs, and waste: it is easy to manufacture prefabricated lumber, from which pieces can be assembled simultaneously (with relatively little labor). This reduces material waste, avoids massive on-site inventory, and minimizes on-site disruption. According to the softwood lumber industry, "Mass timber buildings are roughly 25% faster to construct than concrete buildings and require 90% less construction traffic".
End-of-life
An EPA study showed the typical end-of-life scenario for wood waste from municipal solid waste (MSW), wood packaging, and other miscellaneous wood products in the US. Based on the 2018 data, about 67% of wood waste was landfilled, 16% incinerated with energy recovery, and 17% recycled.
A 2020 study conducted by Edinburgh Napier University demonstrated the proportional waste stream of recovered lumber in the UK. The study showed that timber from municipal solid waste and packaging waste made up 13 and 26% of waste collected. Construction and demolition waste made up the biggest bulk of waste collectively at 52%, with the remaining 10% coming from industry.
In the circular economy
The lumber industry creates a lot of waste, especially in its manufacturing process. From log debarking to finished products, there are several stages of processing that generate a considerable volume of waste, which includes solid wood waste, harmful gases, and residual water.
Wood waste can be recycled at its end of life to make new products. Recycled chips can be used to make wood panels. Such practice reduces the use of virgin raw materials, eliminating emissions that would have otherwise been emitted in its manufacturing.
One of the studies conducted in Hong Kong was done using life-cycle assessment (LCA). The study aimed to assess and compare the environmental impacts of wood waste management from building construction activities using different alternative management scenarios in Hong Kong. Despite various advantages of lumber and its waste, the contribution to the study of the circular economy of lumber is still very small. Some areas where improvements can be made to improve the circularity of lumber is as follows:
First, regulations to support recycled lumber use. For example, establishing grading standards and enforcing penalties for improper disposal, especially in sectors that produce big quantities of wood waste, such as the construction and demolition sector.
Second, creating a stronger supply force. This can be achieved by improving demolition protocol and technology and enhancing the secondary raw materials market through circular business models.
Third, increase demand by introducing incentives to the construction sector and new homeowners to use recycled lumber. This can be in the form of reduced taxes for the construction of the new build.
Secondary raw material
The term secondary raw material denotes waste material that has been recycled and injected back into use as productive material. Lumber has a high potential to be used as a secondary raw material at various stages, as listed below:
Recovery of branches and leaves for use as fertilisers
Timber undergo multiple processing stages before lumber of desired shapes, size, and standards are achieved for commercial use. The process generates a lot of waste which in most cases is disregarded. But being an organic waste, the positive aspect of such waste is that it can be used as a fertiliser or to protect the soil in severe weather conditions.
Recovery of woodchips for thermal energy generation
Waste generated during the manufacturing of lumber products can be used to produce thermal energy. Lumber products after their end-of-life can be downcycled into chips and be used as biomass to produce thermal energy. It is beneficial for industries that need thermal energy.
Circular economy practices offer effective solutions concerning waste. It targets its unnecessary generation through waste reduction, reuse, and recycling. There is no clear explicit evidence of circular economy in the wood panel industry. However, based on the circular economy concept and its characteristics, there are opportunities present in the wood panel industry from the raw material extraction phase to its end-of-life. Therefore, there lies a gap yet to be explored.
| Technology | Building materials | null |
59405 | https://en.wikipedia.org/wiki/Initial%20and%20terminal%20objects | Initial and terminal objects | In category theory, a branch of mathematics, an initial object of a category is an object in such that for every object in , there exists precisely one morphism .
The dual notion is that of a terminal object (also called terminal element): is terminal if for every object in there exists exactly one morphism . Initial objects are also called coterminal or universal, and terminal objects are also called final.
If an object is both initial and terminal, it is called a zero object or null object. A pointed category is one with a zero object.
A strict initial object is one for which every morphism into is an isomorphism.
Examples
The empty set is the unique initial object in Set, the category of sets. Every one-element set (singleton) is a terminal object in this category; there are no zero objects. Similarly, the empty space is the unique initial object in Top, the category of topological spaces and every one-point space is a terminal object in this category.
In the category Rel of sets and relations, the empty set is the unique initial object, the unique terminal object, and hence the unique zero object.
In the category of pointed sets (whose objects are non-empty sets together with a distinguished element; a morphism from to being a function with ), every singleton is a zero object. Similarly, in the category of pointed topological spaces, every singleton is a zero object.
In Grp, the category of groups, any trivial group is a zero object. The trivial object is also a zero object in Ab, the category of abelian groups, Rng the category of pseudo-rings, R-Mod, the category of modules over a ring, and K-Vect, the category of vector spaces over a field. See Zero object (algebra) for details. This is the origin of the term "zero object".
In Ring, the category of rings with unity and unity-preserving morphisms, the ring of integers Z is an initial object. The zero ring consisting only of a single element is a terminal object.
In Rig, the category of rigs with unity and unity-preserving morphisms, the rig of natural numbers N is an initial object. The zero rig, which is the zero ring, consisting only of a single element is a terminal object.
In Field, the category of fields, there are no initial or terminal objects. However, in the subcategory of fields of fixed characteristic, the prime field is an initial object.
Any partially ordered set can be interpreted as a category: the objects are the elements of , and there is a single morphism from to if and only if . This category has an initial object if and only if has a least element; it has a terminal object if and only if has a greatest element.
Cat, the category of small categories with functors as morphisms has the empty category, 0 (with no objects and no morphisms), as initial object and the terminal category, 1 (with a single object with a single identity morphism), as terminal object.
In the category of schemes, Spec(Z), the prime spectrum of the ring of integers, is a terminal object. The empty scheme (equal to the prime spectrum of the zero ring) is an initial object.
A limit of a diagram F may be characterised as a terminal object in the category of cones to F. Likewise, a colimit of F may be characterised as an initial object in the category of co-cones from F.
In the category ChR of chain complexes over a commutative ring R, the zero complex is a zero object.
In a short exact sequence of the form , the initial and terminal objects are the anonymous zero object. This is used frequently in cohomology theories.
Properties
Existence and uniqueness
Initial and terminal objects are not required to exist in a given category. However, if they do exist, they are essentially unique. Specifically, if and are two different initial objects, then there is a unique isomorphism between them. Moreover, if is an initial object then any object isomorphic to is also an initial object. The same is true for terminal objects.
For complete categories there is an existence theorem for initial objects. Specifically, a (locally small) complete category has an initial object if and only if there exist a set ( a proper class) and an -indexed family of objects of such that for any object of , there is at least one morphism for some .
Equivalent formulations
Terminal objects in a category may also be defined as limits of the unique empty diagram . Since the empty category is vacuously a discrete category, a terminal object can be thought of as an empty product (a product is indeed the limit of the discrete diagram , in general). Dually, an initial object is a colimit of the empty diagram and can be thought of as an empty coproduct or categorical sum.
It follows that any functor which preserves limits will take terminal objects to terminal objects, and any functor which preserves colimits will take initial objects to initial objects. For example, the initial object in any concrete category with free objects will be the free object generated by the empty set (since the free functor, being left adjoint to the forgetful functor to Set, preserves colimits).
Initial and terminal objects may also be characterized in terms of universal properties and adjoint functors. Let 1 be the discrete category with a single object (denoted by •), and let be the unique (constant) functor to 1. Then
An initial object in is a universal morphism from • to . The functor which sends • to is left adjoint to U.
A terminal object in is a universal morphism from to •. The functor which sends • to is right adjoint to .
Relation to other categorical constructions
Many natural constructions in category theory can be formulated in terms of finding an initial or terminal object in a suitable category.
A universal morphism from an object to a functor can be defined as an initial object in the comma category . Dually, a universal morphism from to is a terminal object in .
The limit of a diagram is a terminal object in , the category of cones to . Dually, a colimit of is an initial object in the category of cones from .
A representation of a functor to Set is an initial object in the category of elements of .
The notion of final functor (respectively, initial functor) is a generalization of the notion of final object (respectively, initial object).
Other properties
The endomorphism monoid of an initial or terminal object is trivial: .
If a category has a zero object , then for any pair of objects and in , the unique composition is a zero morphism from to .
| Mathematics | Category theory | null |
59407 | https://en.wikipedia.org/wiki/Pea | Pea | Pea (pisum in Latin) is a pulse, vegetable or fodder crop, but the word often refers to the seed or sometimes the pod of this flowering plant species. Carl Linnaeus gave the species the scientific name Pisum sativum in 1753 (meaning cultivated pea). Some sources now treat it as Lathyrus oleraceus; however the need and justification for the change is disputed. Each pod contains several seeds (peas), which can have green or yellow cotyledons when mature. Botanically, pea pods are fruit, since they contain seeds and develop from the ovary of a (pea) flower. The name is also used to describe other edible seeds from the Fabaceae such as the pigeon pea (Cajanus cajan), the cowpea (Vigna unguiculata), the seeds from several species of Lathyrus and is used as a compound form for example Sturt's desert pea.
Peas are annual plants, with a life cycle of one year. They are a cool-season crop grown in many parts of the world; planting can take place from winter to early summer depending on location. The average pea weighs between . The immature peas (and in snow peas the tender pod as well) are used as a vegetable, fresh, frozen or canned; varieties of the species typically called field peas are grown to produce dry peas like the split pea shelled from a matured pod. These are the basis of pease porridge and pea soup, staples of medieval cuisine; in Europe, consuming fresh immature green peas was an innovation of early modern cuisine.
Description
A pea is a most commonly green, occasionally golden yellow, or infrequently purple pod-shaped vegetable, widely grown as a cool-season vegetable crop. The seeds may be planted as soon as the soil temperature reaches , with the plants growing best at temperatures of . They do not thrive in the summer heat of warmer temperate and lowland tropical climates, but do grow well in cooler, high-elevation, tropical areas. Many cultivars reach maturity about 60 days after planting.
Peas have both low-growing and vining cultivars. The vining cultivars grow thin tendrils from leaves that coil around any available support and can climb to be high. A traditional approach to supporting climbing peas is to thrust branches pruned from trees or other woody plants upright into the soil, providing a lattice for the peas to climb. Branches used in this fashion are called pea sticks or sometimes pea brush. Metal fences, twine, or netting supported by a frame are used for the same purpose. In dense plantings, peas give each other some measure of mutual support. Pea plants can self-pollinate.
History
The wild pea is restricted to the Mediterranean Basin and the Near East. The earliest archaeological finds of peas date from the late Neolithic era of current Syria, Anatolia, Israel, Iraq, Jordan and Greece. In Egypt, early finds date from –4400 BC in the Nile delta area, and from c. 3800–3600 BC in Upper Egypt. The pea was also present in Georgia in the 5th millennium BC. Farther east, the finds are younger. Peas were present in Afghanistan c. 2000 BC, in Harappan civilization around modern-day Pakistan and western- and northwestern India in 2250–1750 BC. In the second half of the 2nd millennium BC, this legume crop appears in the Ganges Basin and southern India.
In early times, peas were grown mostly for their dry seeds. From plants growing wild in the Mediterranean Basin, constant selection since the Neolithic dawn of agriculture improved their yield. In the early 3rd century BC, Theophrastus mentions peas among the legumes that are sown late in the winter because of their tenderness. In the first century AD, Columella mentions them in De re rustica, when Roman legionaries still gathered wild peas from the sandy soils of Numidia and Judea to supplement their rations.
In the Middle Ages, field peas are constantly mentioned, as they were the staple that kept famine at bay, as Charles the Good, count of Flanders, noted explicitly in 1124.
Green "garden" peas, eaten immature and fresh, were an innovative luxury of Early Modern Europe. In England, the distinction between field peas and garden peas dates from the early 17th century: John Gerard and John Parkinson both mention garden peas. Sugar peas, which the French called , because they were eaten pods and all, were introduced to France from the market gardens of Holland in the time of Henri IV, through the French ambassador. Green peas were introduced from Genoa to the court of Louis XIV of France in January 1660, with some staged fanfare. A hamper of them was presented before the King. They were shelled by the Savoyan comte de Soissons, who had married a niece of Cardinal Mazarin. Little dishes of peas were then presented to the King, the Queen, Cardinal Mazarin and Monsieur, the king's brother. Immediately established and grown for earliness warmed with manure and protected under glass, they were still a luxurious delicacy in 1696, when Mme de Maintenon and Mme de Sevigné each reported that they were "a fashion, a fury".
The world’s first sweet tasting pea was developed in the 18th century by amateur plant breeder Thomas Edward Knight of Downton, near Salisbury, England.
Modern split peas, with their indigestible skins rubbed off, are a development of the later 19th century.
The top producer of green peas – by far – is China with 12.2 million tons, followed by India (4.8 million tons), USA (0.31 million tons), France (0.23 million tons) and Egypt (0.15 million tons). United Kingdom, Pakistan, Algeria, Peru and Turkey complete the top 10.
Etymology
The term pea originates from the Latin word , which is the latinisation of the Greek (), neuter variant form of () 'pea'. It was adopted into English as the noun pease (plural peasen), as in pease pudding. However, by analogy with other plurals ending in -s, speakers began construing pease as a plural and constructing the singular form by dropping the -s, giving the term pea. This process is known as back-formation.
Composition
Nutrition
Raw green peas are 79% water, 14% carbohydrates, 5% protein, and contain negligible fat (table). In a reference amount of , raw green peas supply of food energy, and are a rich source (20% or more of the Daily Value, DV) of vitamin C (48% DV), vitamin K, thiamine, and manganese, with several B vitamins and dietary minerals in moderate amounts (11–16% DV) (table).
Genome
The pea karyotype consists of seven chromosomes, five of which are acrocentric and two submetacentric. Despite its scientific popularity, its relatively large genome size (4.45Gb) made it challenging to sequence compared to other legumes such as Medicago truncatula and soybeans. The International Pea Genome Sequencing Consortium was formed to develop the first pea reference genome, and the draft assembly was officially announced in September 2019. It covers 88% of the genome (3.92Gb) and predicted 44,791 gene-coding sequences. The pea used for the assembly was the inbred French cultivar "Caméor".
Varieties
Garden peas
There are many varieties (cultivars) of garden peas. Some of the most common varieties are listed here. PMR indicates some degree of powdery mildew resistance; afila types, also called semi-leafless, have clusters of tendrils instead of leaves. Unless otherwise noted these are so called dwarf varieties which grow to an average height of about 1m. Giving the vines support is recommended, but not required. Extra dwarf are suitable for container growing, reaching only about 25 cm. Tall varieties grow to about 2m with support required.
Alaska, 55 days (smooth seeded)
Tom Thumb / Half Pint, 55 days (heirloom, extra dwarf)
Thomas Laxton (heirloom) / Laxton's Progress / Progress #9, 60–65 days
Mr. Big, 60 days, 2000 AAS winner
Little Marvel, 63 days, 1934 AAS winner
Early Perfection, 65 days
Kelvedon Wonder, 65 days, 1997 RHS AGM winner
Sabre, 65 days, PMR
Homesteader / Lincoln, 67 days (heirloom, known as Greenfeast in Australia and New Zealand)
Miragreen, 68 days (tall climber)
Serge, 68 days, PMR, afila
Wando, 68 days
Green Arrow, 70 days
Recruit, 70 days, PMR, afila
Tall Telephone / Alderman, 75 days (heirloom, tall climber)
Edible-pod peas
Some peas lack the tough membrane inside the pod wall and have tender edible pods. There are two main types:
Snow peas have flat pods with thin pod walls. Pods and seeds are eaten when they are very young.
Snap peas or sugar snap peas have rounded pods with thick pod walls. Pods and seeds are eaten before maturity.
The name sugar pea can include both types or be synonymous with either snow peas or snap peas in different dictionaries. Likewise mangetout (; from , 'eat-all pea').
Snow peas and snap peas both belong to Macrocarpon Group, a cultivar group based on the variety Pisum sativum var. macrocarpum Ser. named in 1825. It was described as having very compressed non-leathery edible pods in the original publication.
The scientific name Pisum sativum var. saccharatum Ser. is often misused for snow peas. The variety under this name was described as having sub-leathery and compressed-terete pods and a French name of petit pois. The description is inconsistent with the appearance of snow peas, and therefore botanists have replaced this name with Pisum sativum var. macrocarpum.
Field peas
The field pea is a type of pea sometimes called P. sativum subsp. arvense (L.) Asch. It is also known as dun (grey-brown) pea, Kapucijner pea, or Austrian winter pea, and is one of the oldest domesticated crops, cultivated for at least 7,000 years. Field peas are now grown in many countries for both human consumption and stockfeed. There are several cultivars and colors including blue, dun (brown), maple and white. This pea should not be confused with the cowpea (Vigna unguiculata) which is sometimes called the "field pea" in warmer climates.
It is a climbing annual legume with weak, viny, and relatively succulent stems. Vines often are 4 to 5 feet (120 to 150 cm) long, but when grown alone, field pea's weak stems prevent it from growing more than 1.5 to 2 feet (45 to 60 cm) tall. Leaves have two leaflets and a tendril. Flowers are white, pink, or purple. Pods carry seeds that are large (4,000 seeds/lb), nearly spherical, and white, gray, green, or brown. The root system is relatively shallow and small, but well nodulated.
The field pea is a cool-season legume crop that is grown on over 25 million acres worldwide. It has been an important grain legume crop for millennia, seeds showing domesticated characteristics dating from at least 7000 years ago have been found in archaeological sites around what is now Turkey. Field peas or "dry peas" are marketed as a dry, shelled product for either human or livestock food, unlike the garden pea, which is marketed as a fresh or canned vegetable. The major producing countries of field peas are Russia and China, followed by Canada, Europe, Australia and the United States. Europe, Australia, Canada and the United States raise over 4.5 million acres (18,000 km²) and are major exporters of peas. In 2002, there were approximately 300,000 acres (1,200 km²) of field peas grown in the United States.
Uses
Culinary
In modern times peas are usually boiled or steamed, which breaks down the cell walls and makes them taste sweeter and the nutrients more bioavailable. Along with broad beans and lentils, these formed an important part of the diet of most people in the Middle East, North Africa and Europe during the Middle Ages. By the 17th and 18th centuries, it had become popular to eat peas "green", that is, while they are immature and right after they are picked. New cultivars of peas were developed by the English during this time, which became known as "garden" or "English" peas. The popularity of green peas spread to North America. Thomas Jefferson grew more than 30 cultivars of peas on his estate. With the invention of canning, peas were one of the first vegetables to be canned.
Fresh peas are often eaten boiled and flavored with butter and/or spearmint as a side dish vegetable. Salt and pepper are also commonly added to peas when served. Fresh peas are also used in pot pies, salads and casseroles. Pod peas (snow peas and snap peas) are used in stir-fried dishes, particularly those in American Chinese cuisine. Pea pods do not keep well once picked, and if not used quickly, are best preserved by drying, canning or freezing within a few hours of harvest.
In India, fresh peas are used in various dishes such as aloo matar (curried potatoes with peas) or mattar paneer (paneer cheese with peas), though they can be substituted with frozen peas as well. Peas are also eaten raw, as they are sweet when fresh off the bush. Green peas known as hasiru batani in Kannada are used to make curry and gasi. Split peas are also used to make dal, particularly in Guyana, and Trinidad, where there is a significant population of Indians.
Dried peas are often made into a soup or simply eaten on their own. In Japan, China, Taiwan and some Southeast Asian countries, including Thailand, the Philippines and Malaysia, peas are roasted and salted, and eaten as snacks. In the Philippines, peas, while still in their pods, are a common ingredient in viands and pansit. In the UK, dried yellow or green split peas are used to make pease pudding (or "pease porridge"), a traditional dish. In North America, a similarly traditional dish is split pea soup.
Pea soup is eaten in many other parts of the world, including northern Europe, parts of middle Europe, Russia, Iran, Iraq and India.
In Chinese cuisine, the tender new growth [leaves and stem] (豆苗; ) are commonly used in stir-fries. Much like picking the leaves for tea, the farmers pick the tips off of the pea plant.
In Greece, Tunisia, Turkey, Cyprus, and other parts of the Mediterranean, peas are made into a stew with lamb and potatoes.
In Hungary and Serbia, pea soup is often served with dumplings and spiced with hot paprika.
In the United Kingdom, dried, rehydrated and mashed marrowfat peas, or cooked green split peas, known as mushy peas, are popular, originally in the north of England, but now ubiquitously, and especially as an accompaniment to fish and chips or meat pies, particularly in fish and chip shops. Sodium bicarbonate is sometimes added to soften the peas. In 2005, a poll of 2,000 people revealed the pea to be Britain's seventh favourite culinary vegetable.
Processed peas are mature peas which have been dried, soaked and then heat treated (processed) to prevent spoilage—in the same manner as pasteurizing. Cooked peas are sometimes sold dried and coated with wasabi, salt, or other spices.
In North America pea milk is produced and sold as an alternative to cow milk for a variety of reasons.
Pea sprouts
In East Asia, pea sprouts or shoots (; ) were once dedicated cuisine when the plant was less highly available. Today, when the plant can be easily grown, fresh pea shoots are available in supermarkets or may be grown at home.
Manufacturing
Frozen peas
In order to freeze and preserve peas, they must first be grown, picked, and shelled. Usually, the more tender the peas are, the more likely that they will be used in the final product. The peas must be put through the process of freezing shortly after being picked so that they do not spoil too soon. Once the peas have been selected, they are placed in ice water and allowed to cool. After, they are sprayed with water to remove any residual dirt or dust that may remain on them. The next step is blanching. The peas are boiled for a few minutes to remove any enzymes that may shorten their shelf life. They are then cooled and removed from the water. The final step is the actual freezing to produce the final product. This step may vary considerably; some companies freeze their peas by air blast freezing, where the vegetables are put through a tunnel at high speeds and frozen by cold air. Finally, the peas are packaged and shipped out for retail sale.
Science
In the mid-19th century, Austrian monk Gregor Mendel's observations of pea pods led to the principles of Mendelian genetics, the foundation of modern genetics. He ended up growing and examining about 28,000 pea plants in the course of his experiments.
Mendel chose peas for his experiments because he could grow them easily, pure-bred strains were readily available, and the structure of the flowers protect them from cross-pollination, and cross pollination was easy. Mendel cross-bred tall and dwarf pea plants, green and yellow peas, purple and white flowers, wrinkled and smooth peas, and a few other traits. He then observed the resulting offspring. In each of these cases, one trait is dominant and all the offspring, or Filial-1 (abbreviated F1) generation, showed the dominant trait. Then he allowed the F1 generation to self pollinate and observed their offspring, the Filial-2 (abbreviated F2) generation. The F2 plants had the dominant trait in approximately a 3:1 ratio. He studied later generations of self pollinated plants, and performed crosses to determine the nature of the pollen and egg cells.
Mendel reasoned that each parent had a 'vote' in the appearance of the offspring, and the non-dominant, or recessive, trait appeared only when it was inherited from both parents. He did further experiments that showed each trait is separately inherited. Unwittingly, Mendel had solved a major problem with Charles Darwin's theory of evolution: how new traits were preserved and not blended back into the population, a question Darwin himself did not answer. Mendel's work was published in an obscure Austrian journal and was not rediscovered until about 1900.
Potential for adverse effects
Some people experience allergic reactions to peas, as well as lentils, with vicilin or convicilin as the most common allergens.
Favism, or Fava-bean-ism, is a genetic deficiency of the enzyme glucose-6-phosphate dehydrogenase that affects Jews, other Middle Eastern Semitic peoples, and other descendants of the Mediterranean coastal regions. In this condition, the toxic reaction to eating most, if not all, beans is hemolytic anemia, and in severe cases, the released circulating free hemoglobin causes acute kidney injury.
Nitrogen fixation
Peas, like many legumes, contain symbiotic bacteria called Rhizobia within root nodules of their root systems. These bacteria have the special ability to fix nitrogen from atmospheric, molecular nitrogen () into ammonia (). The chemical reaction is:
Ammonia is then converted to another form, ammonium (), usable by (some) plants by the following reaction:
The root nodules of peas and other legumes are sources of nitrogen that they can use to make amino acids, constituents of proteins. Hence, legumes are good sources of plant protein.
When a pea plant dies in the field, for example following the harvest, all of its remaining nitrogen, incorporated into amino acids inside the remaining plant parts, is released back into the soil. In the soil, the amino acids are converted to nitrate (), that is available to other plants, thereby serving as fertilizer for future crops.
Cultivation
Grading
Pea grading involves sorting peas by size, in which the smallest peas are graded as the highest quality for their tenderness. Brines may be used, in which peas are floated, from which their density can be determined.
Pests and diseases
A variety of diseases affect peas through a number of pathogens, including insects, viruses, bacteria and fungi. In particular, virus disease of peas has worldwide economic importance.
Additionally, insects such as the pea leaf weevil (Sitona lineatus) can damage peas and other pod fruits. The pea leaf weevil is native to Europe, but has spread to other places such as Alberta, Canada. They are about — long and are distinguishable by three light-coloured stripes running length-wise down the thorax. The weevil larvae feed on the root nodules of pea plants, which are essential to the plants' supply of nitrogen, and thus diminish leaf and stem growth. Adult weevils feed on the leaves and create a notched, "c-shaped" appearance on the outside of the leaves.
The pea moth can be a serious pest producing caterpillars the resemble small white maggots in the pea-pods. The caterpillars eat the developing peas making them unsightly and unsuitable for culinary use. Prior to the use of modern insecticides, pea moth caterpillars were a very common sight in pea pods.
| Biology and health sciences | Fabales | null |
59408 | https://en.wikipedia.org/wiki/Pedology | Pedology | Pedology (from Greek: πέδον, pedon, "soil"; and λόγος, logos, "study") is a discipline within soil science which focuses on understanding and characterizing soil formation, evolution, and the theoretical frameworks for modeling soil bodies, often in the context of the natural environment. Pedology is often seen as one of two main branches of soil inquiry, the other being edaphology which is traditionally more agronomically oriented and focuses on how soil properties influence plant communities (natural or cultivated). In studying the fundamental phenomenology of soils, e.g. soil formation (aka pedogenesis), pedologists pay particular attention to observing soil morphology and the geographic distributions of soils, and the placement of soil bodies into larger temporal and spatial contexts. In so doing, pedologists develop systems of soil classification, soil maps, and theories for characterizing temporal and spatial interrelations among soils . There are a few noteworthy sub-disciplines of pedology; namely pedometrics and soil geomorphology. Pedometrics focuses on the development of techniques for quantitative characterization of soils, especially for the purposes of mapping soil properties whereas soil geomorphology studies the interrelationships between geomorphic processes and soil formation.
Overview
Soil is not only a support for vegetation, but it is also the pedosphere, the locus of numerous interactions between climate (water, air, temperature), soil life (micro-organisms, plants, animals) and its residues, the mineral material of the original and added rock, and its position in the landscape. During its formation and genesis, the soil profile slowly deepens and develops characteristic layers, called 'horizons', while a steady state balance is approached.
Soil users (such as agronomists) showed initially little concern in the dynamics of soil. They saw it as medium whose chemical, physical and biological properties were useful for the services of agronomic productivity. On the other hand, pedologists and geologists did not initially focus on the agronomic applications of the soil characteristics (edaphic properties) but upon its relation to the nature and history of landscapes. Today, there is an integration of the two disciplinary approaches as part of landscape and environmental sciences.
Pedologists are now also interested in the practical applications of a good understanding of pedogenesis processes (the evolution and functioning of soils), like interpreting its environmental history and predicting consequences of changes in land use, while agronomists understand that the cultivated soil is a complex medium, often resulting from several thousands of years of evolution. They understand that the current balance is fragile and that only a thorough knowledge of its history makes it possible to ensure its sustainable use.
Concepts
Important pedological concepts include:
Complexity in soil genesis is more common than simplicity.
Soils lie at the interface of Earth's atmosphere, biosphere, hydrosphere and lithosphere. Therefore, a thorough understanding of soils requires some knowledge of meteorology, climatology, ecology, biology, hydrology, geomorphology, geology and many other earth sciences and natural sciences.
Contemporary soils carry imprints of pedogenic processes that were active in the past, although in many cases these imprints are difficult to observe or quantify. Thus, knowledge of paleoecology, palaeogeography, glacial geology and paleoclimatology is important for the recognition and understanding of soil genesis and constitute a basis for predicting future soil changes.
Five major, external factors of formation (climate, organisms, relief, parent material and time), and several smaller, less identifiable ones, drive pedogenic processes and create soil patterns.
Characteristics of soils and soil landscapes, e.g., the number, sizes, shapes and arrangements of soil bodies, each of which is characterized on the basis of soil horizons, degree of internal homogeneity, slope, aspect, landscape position, age and other properties and relationships, can be observed and measured.
Distinctive bioclimatic regimes or combinations of pedogenic processes produce distinctive soils. Thus, distinctive, observable morphological features, e.g., illuvial clay accumulation in B horizons, are produced by certain combinations of pedogenic processes operative over varying periods of time.
Pedogenic (soil-forming) processes act to both create and destroy order (anisotropy) within soils; these processes can proceed simultaneously. The resulting soil profile reflects the balance of these processes, present and past.
The geological Principle of Uniformitarianism applies to soils, i.e., pedogenic processes active in soils today have been operating for long periods of time, back to the time of appearance of organisms on the land surface. These processes do, however, have varying degrees of expression and intensity over space and time.
A succession of different soils may have developed, eroded and/or regressed at any particular site, as soil genetic factors and site factors, e.g., vegetation, sedimentation, geomorphology, change.
There are very few old soils (in a geological sense) because they can be destroyed or buried by geological events, or modified by shifts in climate by virtue of their vulnerable position at the surface of the earth. Little of the soil continuum dates back beyond the Tertiary period and most soils and land surfaces are no older than the Pleistocene Epoch. However, preserved/lithified soils (paleosols) are an almost ubiquitous feature in terrestrial (land-based) environments throughout most of geologic time. Since they record evidence of ancient climate change, they present immense utility in understanding climate evolution throughout geologic history.
Knowledge and understanding of the genesis of a soil is important in its classification and mapping.
Soil classification systems cannot be based entirely on perceptions of genesis, however, because genetic processes are seldom observed and because pedogenic processes change over time.
Knowledge of soil genesis is imperative and basic to soil use and management. Human influence on, or adjustment to, the factors and processes of soil formation can be best controlled and planned using knowledge about soil genesis.
Soils are natural clay factories (clay includes both clay mineral structures and particles less than 2 μm in diameter). Shales worldwide are, to a considerable extent, simply soil clays that have been formed in the pedosphere and eroded and deposited in the ocean basins, to become lithified at a later date.
Notable pedologists
Olivier de Serres
Vasily V. Dokuchaev
Friedrich Albert Fallou
Konstantin D. Glinka
Eugene W. Hilgard
Francis D. Hole
Hans Jenny
Curtis F. Marbut
Bernard Palissy
| Physical sciences | Soil science | Earth science |
59413 | https://en.wikipedia.org/wiki/Soil%20science | Soil science | Soil science is the study of soil as a natural resource on the surface of the Earth including soil formation, classification and mapping; physical, chemical, biological, and fertility properties of soils; and these properties in relation to the use and management of soils.
The main branches of soil science are pedology ― the study of formation, chemistry, morphology, and classification of soil ― and edaphology ― the study of how soils interact with living things, especially plants. Sometimes terms which refer to those branches are used as if synonymous with soil science. The diversity of names associated with this discipline is related to the various associations concerned. Indeed, engineers, agronomists, chemists, geologists, physical geographers, ecologists, biologists, microbiologists, silviculturists, sanitarians, archaeologists, and specialists in regional planning, all contribute to further knowledge of soils and the advancement of the soil sciences.
Soil scientists have raised concerns about how to preserve soil and arable land in a world with a growing population, possible future water crisis, increasing per capita food consumption, and land degradation.
Fields of study
Soil occupies the pedosphere, one of Earth's spheres that the geosciences use to organize the Earth conceptually. This is the conceptual perspective of pedology and edaphology, the two main branches of soil science. Pedology is the study of soil in its natural setting. Edaphology is the study of soil in relation to soil-dependent uses. Both branches apply a combination of soil physics, soil chemistry, and soil biology. Due to the numerous interactions between the biosphere, atmosphere and hydrosphere that are hosted within the pedosphere, more integrated, less soil-centric concepts are also valuable. Many concepts essential to understanding soil come from individuals not identifiable strictly as soil scientists. This highlights the interdisciplinary nature of soil concepts.
Research
Exploring the diversity and dynamics of soil continues to yield fresh discoveries and insights. New avenues of soil research are compelled by a need to understand soil in the context of climate change, greenhouse gases, and carbon sequestration. Interest in maintaining the planet's biodiversity and in exploring past cultures has also stimulated renewed interest in achieving a more refined understanding of soil.
Mapping
Classification
In 1998, the World Reference Base for Soil Resources (WRB) replaced the FAO soil classification as the international soil classification system. The currently valid version of WRB is the 4th edition, 2022. The FAO soil classification, in turn, borrowed from modern soil classification concepts, including USDA soil taxonomy.
WRB is based mainly on soil morphology as an expression of pedogenesis. A major difference with USDA soil taxonomy is that soil climate is not part of the system, except insofar as climate influences soil profile characteristics.
Many other classification schemes exist, including vernacular systems. The structure in vernacular systems is either nominal (giving unique names to soils or landscapes) or descriptive (naming soils by their characteristics such as red, hot, fat, or sandy). Soils are distinguished by obvious characteristics, such as physical appearance (e.g., color, texture, landscape position), performance (e.g., production capability, flooding), and accompanying vegetation. A vernacular distinction familiar to many is classifying texture as heavy or light. Light soil content and better structure take less effort to turn and cultivate. Light soils do not necessarily weigh less than heavy soils on an air dry basis, nor do they have more porosity.
History
The earliest known soil classification system comes from China, appearing in the book Yu Gong (5th century BCE), where the soil was divided into three categories and nine classes, depending on its color, texture and hydrology.
Contemporaries Friedrich Albert Fallou (the German founder of modern soil science) and Vasily Dokuchaev (the Russian founder of modern soil science) are both credited with being among the first to identify soil as a resource whose distinctness and complexity deserved to be separated conceptually from geology and crop production and treated as a whole. As a founding father of soil science, Fallou has primacy in time. Fallou was working on the origins of soil before Dokuchaev was born; however Dokuchaev's work was more extensive and is considered to be the more significant to modern soil theory than Fallou's.
Previously, soil had been considered a product of chemical transformations of rocks, a dead substrate from which plants derive nutritious elements. Soil and bedrock were in fact equated. Dokuchaev considers the soil as a natural body having its own genesis and its own history of development, a body with complex and multiform processes taking place within it. The soil is considered as different from bedrock. The latter becomes soil under the influence of a series of soil-formation factors (climate, vegetation, country, relief and age). According to him, soil should be called the "daily" or outward horizons of rocks regardless of the type; they are changed naturally by the common effect of water, air and various kinds of living and dead organisms.
A 1914 encyclopedic definition: "the different forms of earth on the surface of the rocks, formed by the breaking down or weathering of rocks". serves to illustrate the historic view of soil which persisted from the 19th century. Dokuchaev's late 19th century soil concept developed in the 20th century to one of soil as earthy material that has been altered by living processes. A corollary concept is that soil without a living component is simply a part of Earth's outer layer.
Further refinement of the soil concept is occurring in view of an appreciation of energy transport and transformation within soil. The term is popularly applied to the material on the surface of the Earth's moon and Mars, a usage acceptable within a portion of the scientific community. Accurate to this modern understanding of soil is Nikiforoff's 1959 definition of soil as the "excited skin of the sub aerial part of the Earth's crust".
Areas of practice
Academically, soil scientists tend to be drawn to one of five areas of specialization: microbiology, pedology, edaphology, physics, or chemistry. Yet the work specifics are very much dictated by the challenges facing our civilization's desire to sustain the land that supports it, and the distinctions between the sub-disciplines of soil science often blur in the process. Soil science professionals commonly stay current in soil chemistry, soil physics, soil microbiology, pedology, and applied soil science in related disciplines.
One exciting effort drawing in soil scientists in the U.S. is the Soil Quality Initiative. Central to the Soil Quality Initiative is developing indices of soil health and then monitoring them in a way that gives us long-term (decade-to-decade) feedback on our performance as stewards of the planet. The effort includes understanding the functions of soil microbiotic crusts and exploring the potential to sequester atmospheric carbon in soil organic matter. Relating the concept of agriculture to soil quality, however, has not been without its share of controversy and criticism, including critiques by Nobel Laureate Norman Borlaug and World Food Prize Winner Pedro Sanchez.
A more traditional role for soil scientists has been to map soils. Almost every area in the United States now has a published soil survey, including interpretive tables on how soil properties support or limit activities and uses. An internationally accepted soil taxonomy allows uniform communication of soil characteristics and soil functions. National and international soil survey efforts have given the profession unique insights into landscape-scale functions. The landscape functions that soil scientists are called upon to address in the field seem to fall roughly into six areas:
Land-based treatment of wastes
Septic system
Manure
Municipal biosolids
Food and fiber processing waste
Identification and protection of environmentally critical areas
Sensitive and unstable soils
Wetlands
Unique soil situations that support valuable habitat, and ecosystem diversity
Management for optimum land productivity
Silviculture
Agronomy
Nutrient management
Water management
Native vegetation
Grazing
Management for optimum water quality
Stormwater management
Sediment and erosion control
Remediation and restoration of damaged lands
Mine reclamation
Flood and storm damage
Contamination
Sustainability of desired uses
Soil conservation
There are also practical applications of soil science that might not be apparent from looking at a published soil survey.
Radiometric dating: specifically a knowledge of local pedology is used to date prior activity at the site
Stratification (archeology) where soil formation processes and preservative qualities can inform the study of archaeological sites
Geological phenomena
Landslides
Active faults
Altering soils to achieve new uses
Vitrification to contain radioactive wastes
Enhancing soil microbial capabilities in degrading contaminants (bioremediation).
Carbon sequestration
Environmental soil science
Pedology
Soil genesis
Pedometrics
Soil morphology
Soil micromorphology
Soil classification
USDA soil taxonomy
World Reference Base for Soil Resources
Soil biology
Soil microbiology
Soil chemistry
Soil biochemistry
Soil mineralogy
Soil physics
Pedotransfer function
Soil mechanics and engineering
Soil hydrology, hydropedology
Fields of application in soil science
Climate change
Ecosystem studies
Pedotransfer function
Soil fertility / Nutrient management
Soil management
Soil survey
Standard methods of analysis
Watershed and wetland studies
Land Suitability classification
Related disciplines
Agricultural sciences
Agricultural soil science
Agrophysics science
Irrigation management
Anthropology
archaeological stratigraphy
Environmental science
Landscape ecology
Physical geography
Geomorphology
Geology
Biogeochemistry
Geomicrobiology
Hydrology
Hydrogeology
Waste management
Wetland science
Depression storage capacity
Depression storage capacity, in soil science, is the ability of a particular area of land to retain water in its pits and depressions, thus preventing it from flowing. Depression storage capacity, along with infiltration capacity, is one of the main factors involved in Horton overland flow, whereby water volume surpasses both infiltration and depression storage capacity and begins to flow horizontally across land, possibly leading to flooding and soil erosion. The study of land's depression storage capacity is important in the fields of geology, ecology, and especially hydrology.
| Physical sciences | Pedology | null |
59414 | https://en.wikipedia.org/wiki/Nitrogen%20cycle | Nitrogen cycle | The nitrogen cycle is the biogeochemical cycle by which nitrogen is converted into multiple chemical forms as it circulates among atmospheric, terrestrial, and marine ecosystems. The conversion of nitrogen can be carried out through both biological and physical processes. Important processes in the nitrogen cycle include fixation, ammonification, nitrification, and denitrification. The majority of Earth's atmosphere (78%) is atmospheric nitrogen, making it the largest source of nitrogen. However, atmospheric nitrogen has limited availability for biological use, leading to a scarcity of usable nitrogen in many types of ecosystems.
The nitrogen cycle is of particular interest to ecologists because nitrogen availability can affect the rate of key ecosystem processes, including primary production and decomposition. Human activities such as fossil fuel combustion, use of artificial nitrogen fertilizers, and release of nitrogen in wastewater have dramatically altered the global nitrogen cycle. Human modification of the global nitrogen cycle can negatively affect the natural environment system and also human health.
Processes
Nitrogen is present in the environment in a wide variety of chemical forms including organic nitrogen, ammonium (), nitrite (), nitrate (), nitrous oxide (), nitric oxide (NO) or inorganic nitrogen gas (). Organic nitrogen may be in the form of a living organism, humus or in the intermediate products of organic matter decomposition. The processes in the nitrogen cycle is to transform nitrogen from one form to another. Many of those processes are carried out by microbes, either in their effort to harvest energy or to accumulate nitrogen in a form needed for their growth. For example, the nitrogenous wastes in animal urine are broken down by nitrifying bacteria in the soil to be used by plants. The diagram alongside shows how these processes fit together to form the nitrogen cycle.
Nitrogen fixation
The conversion of nitrogen gas () into nitrates and nitrites through atmospheric, industrial and biological processes is called nitrogen fixation. Atmospheric nitrogen must be processed, or "fixed", into a usable form to be taken up by plants. Between 5 and 10 billion kg per year are fixed by lightning strikes, but most fixation is done by free-living or symbiotic bacteria known as diazotrophs. These bacteria have the nitrogenase enzyme that combines gaseous nitrogen with hydrogen to produce ammonia, which is converted by the bacteria into other organic compounds. Most biological nitrogen fixation occurs by the activity of molybdenum (Mo)-nitrogenase, found in a wide variety of bacteria and some Archaea. Mo-nitrogenase is a complex two-component enzyme that has multiple metal-containing prosthetic groups. An example of free-living bacteria is Azotobacter. Symbiotic nitrogen-fixing bacteria such as Rhizobium usually live in the root nodules of legumes (such as peas, alfalfa, and locust trees). Here they form a mutualistic relationship with the plant, producing ammonia in exchange for carbohydrates. Because of this relationship, legumes will often increase the nitrogen content of nitrogen-poor soils. A few non-legumes can also form such symbioses. Today, about 30% of the total fixed nitrogen is produced industrially using the Haber-Bosch process, which uses high temperatures and pressures to convert nitrogen gas and a hydrogen source (natural gas or petroleum) into ammonia.
Assimilation
Plants can absorb nitrate or ammonium from the soil by their root hairs. If nitrate is absorbed, it is first reduced to nitrite ions and then ammonium ions for incorporation into amino acids, nucleic acids, and chlorophyll. In plants that have a symbiotic relationship with rhizobia, some nitrogen is assimilated in the form of ammonium ions directly from the nodules. It is now known that there is a more complex cycling of amino acids between Rhizobia bacteroids and plants. The plant provides amino acids to the bacteroids so ammonia assimilation is not required and the bacteroids pass amino acids (with the newly fixed nitrogen) back to the plant, thus forming an interdependent relationship. While many animals, fungi, and other heterotrophic organisms obtain nitrogen by ingestion of amino acids, nucleotides, and other small organic molecules, other heterotrophs (including many bacteria) are able to utilize inorganic compounds, such as ammonium as sole N sources. Utilization of various N sources is carefully regulated in all organisms.
Ammonification
When a plant or animal dies or an animal expels waste, the initial form of nitrogen is organic. Bacteria or fungi convert the organic nitrogen within the remains back into ammonium (), a process called ammonification or mineralization. Enzymes involved are:
GS: Gln Synthetase (cytosolic & plastic)
GOGAT: Glu 2-oxoglutarate aminotransferase (Ferredoxin & NADH-dependent)
GDH: Glu Dehydrogenase:
Minor role in ammonium assimilation.
Important in amino acid catabolism.
Nitrification
The conversion of ammonium to nitrate is performed primarily by soil-living bacteria and other nitrifying bacteria. In the primary stage of nitrification, the oxidation of ammonium () is performed by bacteria such as the Nitrosomonas species, which converts ammonia to nitrites (). Other bacterial species such as Nitrobacter, are responsible for the oxidation of the nitrites () into nitrates (). It is important for the ammonia () to be converted to nitrates or nitrites because ammonia gas is toxic to plants.
Due to their very high solubility and because soils are highly unable to retain anions, nitrates can enter groundwater. Elevated nitrate in groundwater is a concern for drinking water use because nitrate can interfere with blood-oxygen levels in infants and cause methemoglobinemia or blue-baby syndrome. Where groundwater recharges stream flow, nitrate-enriched groundwater can contribute to eutrophication, a process that leads to high algal population and growth, especially blue-green algal populations. While not directly toxic to fish life, like ammonia, nitrate can have indirect effects on fish if it contributes to this eutrophication. Nitrogen has contributed to severe eutrophication problems in some water bodies. Since 2006, the application of nitrogen fertilizer has been increasingly controlled in Britain and the United States. This is occurring along the same lines as control of phosphorus fertilizer, restriction of which is normally considered essential to the recovery of eutrophied waterbodies.
Denitrification
Denitrification is the reduction of nitrates back into nitrogen gas (), completing the nitrogen cycle. This process is performed by bacterial species such as Pseudomonas and Paracoccus, under anaerobic conditions. They use the nitrate as an electron acceptor in the place of oxygen during respiration. These facultatively (meaning optionally) anaerobic bacteria can also live in aerobic conditions. Denitrification happens in anaerobic conditions e.g. waterlogged soils. The denitrifying bacteria use nitrates in the soil to carry out respiration and consequently produce nitrogen gas, which is inert and unavailable to plants. Denitrification occurs in free-living microorganisms as well as obligate symbionts of anaerobic ciliates.
Dissimilatory nitrate reduction to ammonium
Dissimilatory nitrate reduction to ammonium (DNRA), or nitrate/nitrite ammonification, is an anaerobic respiration process. Microbes which undertake DNRA oxidise organic matter and use nitrate as an electron acceptor, reducing it to nitrite, then ammonium (). Both denitrifying and nitrate ammonification bacteria will be competing for nitrate in the environment, although DNRA acts to conserve bioavailable nitrogen as soluble ammonium rather than producing dinitrogen gas.
Anaerobic ammonia oxidation
The ANaerobic AMMonia OXidation process is also known as the ANAMMOX process, an abbreviation coined by joining the first syllables of each of these three words. This biological process is a redox comproportionation reaction, in which ammonia (the reducing agent giving electrons) and nitrite (the oxidizing agent accepting electrons) transfer three electrons and are converted into one molecule of diatomic nitrogen () gas and two water molecules. This process makes up a major proportion of nitrogen conversion in the oceans. The stoichiometrically balanced formula for the ANAMMOX chemical reaction can be written as following, where an ammonium ion includes the ammonia molecule, its conjugated base:
(ΔG° = ).
This an exergonic process (here also an exothermic reaction) releasing energy, as indicated by the negative value of ΔG°, the difference in Gibbs free energy between the products of reaction and the reagents.
Other processes
Though nitrogen fixation is the primary source of plant-available nitrogen in most ecosystems, in areas with nitrogen-rich bedrock, the breakdown of this rock also serves as a nitrogen source. Nitrate reduction is also part of the iron cycle, under anoxic conditions Fe(II) can donate an electron to and is oxidized to Fe(III) while is reduced to , and depending on the conditions and microbial species involved. The fecal plumes of cetaceans also act as a junction in the marine nitrogen cycle, concentrating nitrogen in the epipelagic zones of ocean environments before its dispersion through various marine layers, ultimately enhancing oceanic primary productivity.
Marine nitrogen cycle
The nitrogen cycle is an important process in the ocean as well. While the overall cycle is similar, there are different players and modes of transfer for nitrogen in the ocean. Nitrogen enters the water through the precipitation, runoff, or as from the atmosphere. Nitrogen cannot be utilized by phytoplankton as so it must undergo nitrogen fixation which is performed predominately by cyanobacteria. Without supplies of fixed nitrogen entering the marine cycle, the fixed nitrogen would be used up in about 2000 years. Phytoplankton need nitrogen in biologically available forms for the initial synthesis of organic matter. Ammonia and urea are released into the water by excretion from plankton. Nitrogen sources are removed from the euphotic zone by the downward movement of the organic matter. This can occur from sinking of phytoplankton, vertical mixing, or sinking of waste of vertical migrators. The sinking results in ammonia being introduced at lower depths below the euphotic zone. Bacteria are able to convert ammonia to nitrite and nitrate but they are inhibited by light so this must occur below the euphotic zone. Ammonification or Mineralization is performed by bacteria to convert organic nitrogen to ammonia. Nitrification can then occur to convert the ammonium to nitrite and nitrate. Nitrate can be returned to the euphotic zone by vertical mixing and upwelling where it can be taken up by phytoplankton to continue the cycle. can be returned to the atmosphere through denitrification.
Ammonium is thought to be the preferred source of fixed nitrogen for phytoplankton because its assimilation does not involve a redox reaction and therefore requires little energy. Nitrate requires a redox reaction for assimilation but is more abundant so most phytoplankton have adapted to have the enzymes necessary to undertake this reduction (nitrate reductase). There are a few notable and well-known exceptions that include most Prochlorococcus and some Synechococcus that can only take up nitrogen as ammonium.
The nutrients in the ocean are not uniformly distributed. Areas of upwelling provide supplies of nitrogen from below the euphotic zone. Coastal zones provide nitrogen from runoff and upwelling occurs readily along the coast. However, the rate at which nitrogen can be taken up by phytoplankton is decreased in oligotrophic waters year-round and temperate water in the summer resulting in lower primary production. The distribution of the different forms of nitrogen varies throughout the oceans as well.
Nitrate is depleted in near-surface water except in upwelling regions. Coastal upwelling regions usually have high nitrate and chlorophyll levels as a result of the increased production. However, there are regions of high surface nitrate but low chlorophyll that are referred to as HNLC (high nitrogen, low chlorophyll) regions. The best explanation for HNLC regions relates to iron scarcity in the ocean, which may play an important part in ocean dynamics and nutrient cycles. The input of iron varies by region and is delivered to the ocean by dust (from dust storms) and leached out of rocks. Iron is under consideration as the true limiting element to ecosystem productivity in the ocean.
Ammonium and nitrite show a maximum concentration at 50–80 m (lower end of the euphotic zone) with decreasing concentration below that depth. This distribution can be accounted for by the fact that nitrite and ammonium are intermediate species. They are both rapidly produced and consumed through the water column. The amount of ammonium in the ocean is about 3 orders of magnitude less than nitrate. Between ammonium, nitrite, and nitrate, nitrite has the fastest turnover rate. It can be produced during nitrate assimilation, nitrification, and denitrification; however, it is immediately consumed again.
New vs. regenerated nitrogen
Nitrogen entering the euphotic zone is referred to as new nitrogen because it is newly arrived from outside the productive layer. The new nitrogen can come from below the euphotic zone or from outside sources. Outside sources are upwelling from deep water and nitrogen fixation. If the organic matter is eaten, respired, delivered to the water as ammonia, and re-incorporated into organic matter by phytoplankton it is considered recycled/regenerated production.
New production is an important component of the marine environment. One reason is that only continual input of new nitrogen can determine the total capacity of the ocean to produce a sustainable fish harvest. Harvesting fish from regenerated nitrogen areas will lead to a decrease in nitrogen and therefore a decrease in primary production. This will have a negative effect on the system. However, if fish are harvested from areas of new nitrogen the nitrogen will be replenished.
Future acidification
As illustrated by the diagram on the right, additional carbon dioxide () is absorbed by the ocean and reacts with water, carbonic acid () is formed and broken down into both bicarbonate () and hydrogen () ions (gray arrow), which reduces bioavailable carbonate () and decreases ocean pH (black arrow). This is likely to enhance nitrogen fixation by diazotrophs (gray arrow), which utilize ions to convert nitrogen into bioavailable forms such as ammonia () and ammonium ions (). However, as pH decreases, and more ammonia is converted to ammonium ions (gray arrow), there is less oxidation of ammonia to nitrite (NO), resulting in an overall decrease in nitrification and denitrification (black arrows). This in turn would lead to a further build-up of fixed nitrogen in the ocean, with the potential consequence of eutrophication. Gray arrows represent an increase while black arrows represent a decrease in the associated process.
Human influences on the nitrogen cycle
As a result of extensive cultivation of legumes (particularly soy, alfalfa, and clover), growing use of the Haber–Bosch process in the production of chemical fertilizers, and pollution emitted by vehicles and industrial plants, human beings have more than doubled the annual transfer of nitrogen into biologically available forms. In addition, humans have significantly contributed to the transfer of nitrogen trace gases from Earth to the atmosphere and from the land to aquatic systems. Human alterations to the global nitrogen cycle are most intense in developed countries and in Asia, where vehicle emissions and industrial agriculture are highest.
Generation of Nr, reactive nitrogen, has increased over 10 fold in the past century due to global industrialisation. This form of nitrogen follows a cascade through the biosphere via a variety of mechanisms, and is accumulating as the rate of its generation is greater than the rate of denitrification.
Nitrous oxide () has risen in the atmosphere as a result of agricultural fertilization, biomass burning, cattle and feedlots, and industrial sources. has deleterious effects in the stratosphere, where it breaks down and acts as a catalyst in the destruction of atmospheric ozone. Nitrous oxide is also a greenhouse gas and is currently the third largest contributor to global warming, after carbon dioxide and methane. While not as abundant in the atmosphere as carbon dioxide, it is, for an equivalent mass, nearly 300 times more potent in its ability to warm the planet.
Ammonia () in the atmosphere has tripled as the result of human activities. It is a reactant in the atmosphere, where it acts as an aerosol, decreasing air quality and clinging to water droplets, eventually resulting in nitric acid (HNO3) that produces acid rain. Atmospheric ammonia and nitric acid also damage respiratory systems.
The very high temperature of lightning naturally produces small amounts of , , and , but high-temperature combustion has contributed to a 6- or 7-fold increase in the flux of to the atmosphere. Its production is a function of combustion temperature - the higher the temperature, the more is produced. Fossil fuel combustion is a primary contributor, but so are biofuels and even the burning of hydrogen. However, the rate that hydrogen is directly injected into the combustion chambers of internal combustion engines can be controlled to prevent the higher combustion temperatures that produce .
Ammonia and nitrous oxides actively alter atmospheric chemistry. They are precursors of tropospheric (lower atmosphere) ozone production, which contributes to smog and acid rain, damages plants and increases nitrogen inputs to ecosystems. Ecosystem processes can increase with nitrogen fertilization, but anthropogenic input can also result in nitrogen saturation, which weakens productivity and can damage the health of plants, animals, fish, and humans.
Decreases in biodiversity can also result if higher nitrogen availability increases nitrogen-demanding grasses, causing a degradation of nitrogen-poor, species-diverse heathlands.
Consequence of human modification of the nitrogen cycle
Impacts on natural systems
Increasing levels of nitrogen deposition is shown to have several adverse effects on both terrestrial and aquatic ecosystems. Nitrogen gases and aerosols can be directly toxic to certain plant species, affecting the aboveground physiology and growth of plants near large point sources of nitrogen pollution. Changes to plant species may also occur as nitrogen compound accumulation increases availability in a given ecosystem, eventually changing the species composition, plant diversity, and nitrogen cycling. Ammonia and ammonium – two reduced forms of nitrogen – can be detrimental over time due to increased toxicity toward sensitive species of plants, particularly those that are accustomed to using nitrate as their source of nitrogen, causing poor development of their roots and shoots. Increased nitrogen deposition also leads to soil acidification, which increases base cation leaching in the soil and amounts of aluminum and other potentially toxic metals, along with decreasing the amount of nitrification occurring and increasing plant-derived litter. Due to the ongoing changes caused by high nitrogen deposition, an environment's susceptibility to ecological stress and disturbance – such as pests and pathogens – may increase, thus making it less resilient to situations that otherwise would have little impact on its long-term vitality.
Additional risks posed by increased availability of inorganic nitrogen in aquatic ecosystems include water acidification; eutrophication of fresh and saltwater systems; and toxicity issues for animals, including humans. Eutrophication often leads to lower dissolved oxygen levels in the water column, including hypoxic and anoxic conditions, which can cause death of aquatic fauna. Relatively sessile benthos, or bottom-dwelling creatures, are particularly vulnerable because of their lack of mobility, though large fish kills are not uncommon. Oceanic dead zones near the mouth of the Mississippi in the Gulf of Mexico are a well-known example of algal bloom-induced hypoxia. The New York Adirondack Lakes, Catskills, Hudson Highlands, Rensselaer Plateau and parts of Long Island display the impact of nitric acid rain deposition, resulting in the killing of fish and many other aquatic species.
Ammonia () is highly toxic to fish, and the level of ammonia discharged from wastewater treatment facilities must be closely monitored. Nitrification via aeration before discharge is often desirable to prevent fish deaths. Land application can be an attractive alternative to aeration.
Impacts on human health: nitrate accumulation in drinking water
Leakage of Nr (reactive nitrogen) from human activities can cause nitrate accumulation in the natural water environment, which can create harmful impacts on human health. Excessive use of N-fertilizer in agriculture has been a significant source of nitrate pollution in groundwater and surface water. Due to its high solubility and low retention by soil, nitrate can easily escape from the subsoil layer to the groundwater, causing nitrate pollution. Some other non-point sources for nitrate pollution in groundwater originate from livestock feeding, animal and human contamination, and municipal and industrial waste. Since groundwater often serves as the primary domestic water supply, nitrate pollution can be extended from groundwater to surface and drinking water during potable water production, especially for small community water supplies, where poorly regulated and unsanitary waters are used.
The WHO standard for drinking water is 50 mg L−1 for short-term exposure, and for 3 mg L−1 chronic effects. Once it enters the human body, nitrate can react with organic compounds through nitrosation reactions in the stomach to form nitrosamines and nitrosamides, which are involved in some types of cancers (e.g., oral cancer and gastric cancer).
Impacts on human health: air quality
Human activities have also dramatically altered the global nitrogen cycle by producing nitrogenous gases associated with global atmospheric nitrogen pollution. There are multiple sources of atmospheric reactive nitrogen (Nr) fluxes. Agricultural sources of reactive nitrogen can produce atmospheric emission of ammonia (), nitrogen oxides () and nitrous oxide (). Combustion processes in energy production, transportation, and industry can also form new reactive nitrogen via the emission of , an unintentional waste product. When those reactive nitrogens are released into the lower atmosphere, they can induce the formation of smog, particulate matter, and aerosols, all of which are major contributors to adverse health effects on human health from air pollution. In the atmosphere, can be oxidized to nitric acid (), and it can further react with to form ammonium nitrate (), which facilitates the formation of particulate nitrate. Moreover, can react with other acid gases (sulfuric and hydrochloric acids) to form ammonium-containing particles, which are the precursors for the secondary organic aerosol particles in photochemical smog.
| Physical sciences | Earth science basics: General | Earth science |
59416 | https://en.wikipedia.org/wiki/Soil%20erosion | Soil erosion | Soil erosion is the denudation or wearing away of the upper layer of soil. It is a form of soil degradation. This natural process is caused by the dynamic activity of erosive agents, that is, water, ice (glaciers), snow, air (wind), plants, and animals (including humans). In accordance with these agents, erosion is sometimes divided into water erosion, glacial erosion, snow erosion, wind (aeolian) erosion, zoogenic erosion and anthropogenic erosion such as tillage erosion.
Soil erosion may be a slow process that continues relatively unnoticed, or it may occur at an alarming rate causing a serious loss of topsoil. The loss of soil from farmland may be reflected in reduced crop production potential, lower surface water quality and damaged drainage networks. Soil erosion could also cause sinkholes.
Human activities have increased by 10–50 times the rate at which erosion is occurring world-wide.
Excessive (or accelerated) erosion causes both "on-site" and "off-site" problems. On-site impacts include decreases in agricultural productivity and (on natural landscapes) ecological collapse, both because of loss of the nutrient-rich upper soil layers. In some cases, the eventual result is desertification. Off-site effects include sedimentation of waterways and eutrophication of water bodies, as well as sediment-related damage to roads and houses. Water and wind erosion are the two primary causes of land degradation; combined, they are responsible for about 84% of the global extent of degraded land, making excessive erosion one of the most significant environmental problems worldwide.
Intensive agriculture, deforestation, roads, acid rains, anthropogenic climate change and urban sprawl are amongst the most significant human activities in regard to their effect on stimulating erosion. However, there are many prevention and remediation practices that can curtail or limit erosion of vulnerable soils.
Physical processes
Rainfall and surface runoff
Rainfall, and the surface runoff which may result from rainfall, produces four main types of soil erosion: splash erosion, sheet erosion, rill erosion, and gully erosion. Splash erosion is generally seen as the first and least severe stage in the soil erosion process, which is followed by sheet erosion, then rill erosion and finally gully erosion (the most severe of the four).
In splash erosion, the impact of a falling raindrop creates a small crater in the soil, ejecting soil particles. The distance these soil particles travel can be as much as 0.6 m (two feet) vertically and 1.5 m (five feet) horizontally on level ground.
If the soil is saturated, or if the rainfall rate is greater than the rate at which water can infiltrate into the soil, surface runoff occurs. If the runoff has sufficient flow energy, it will transport loosened soil particles (sediment) down the slope. Sheet erosion is the transport of loosened soil particles by overland flow.
Rill erosion refers to the development of small, ephemeral concentrated flow paths which function as both sediment source and sediment delivery systems for erosion on hillslopes. Generally, where water erosion rates on disturbed upland areas are greatest, rills are active. Flow depths in rills are typically of the order of a few centimeters (about an inch) or less and along-channel slopes may be quite steep. This means that rills exhibit hydraulic physics very different from water flowing through the deeper wider channels of streams and rivers.
Gully erosion occurs when runoff water accumulates and rapidly flows in narrow channels during or immediately after heavy rains or melting snow, removing soil to a considerable depth. Another cause of gully erosion is grazing, which often results in ground compaction. Because the soil is exposed, it loses the ability to absorb excess water, and erosion can develop in susceptible areas.
Rivers and streams
Valley or stream erosion occurs with continued water flow along a linear feature. The erosion is both downward, deepening the valley, and headward, extending the valley into the hillside, creating head cuts and steep banks. In the earliest stage of stream erosion, the erosive activity is dominantly vertical, the valleys have a typical V cross-section and the stream gradient is relatively steep. When some base level is reached, the erosive activity switches to lateral erosion, which widens the valley floor and creates a narrow floodplain. The stream gradient becomes nearly flat, and lateral deposition of sediments becomes important as the stream meanders across the valley floor. In all stages of stream erosion, by far the most erosion occurs during times of flood, when more and faster-moving water is available to carry a larger sediment load. In such processes, it is not the water alone that erodes: suspended abrasive particles, pebbles and boulders can also act erosively as they traverse a surface, in a process known as traction.
Bank erosion is the wearing away of the banks of a stream or river. This is distinguished from changes on the bed of the watercourse, which is referred to as scour. Erosion and changes in the form of river banks may be measured by inserting metal rods into the bank and marking the position of the bank surface along the rods at different times.
Thermal erosion is the result of melting and weakening permafrost due to moving water. It can occur both along rivers and at the coast. Rapid river channel migration observed in the Lena River of Siberia is due to thermal erosion, as these portions of the banks are composed of permafrost-cemented non-cohesive materials. Much of this erosion occurs as the weakened banks fail in large slumps. Thermal erosion also affects the Arctic coast, where wave action and near-shore temperatures combine to undercut permafrost bluffs along the shoreline and cause them to fail. Annual erosion rates along a segment of the Beaufort Sea shoreline averaged per year from 1955 to 2002.
Floods
At extremely high flows, kolks, or vortices are formed by large volumes of rapidly rushing water. Kolks cause extreme local erosion, plucking bedrock and creating pothole-type geographical features called rock-cut basins. Examples can be seen in the flood regions result from glacial Lake Missoula, which created the channeled scablands in the Columbia Basin region of eastern Washington.
Wind erosion
Wind erosion is a major geomorphological force, especially in arid and semi-arid regions. It is also a major source of land degradation, evaporation, desertification, harmful airborne dust, and crop damage—especially after being increased far above natural rates by human activities such as deforestation, urbanization, and agriculture.
Wind erosion is of two primary varieties: deflation, where the wind picks up and carries away loose particles; and abrasion, where surfaces are worn down as they are struck by airborne particles carried by wind. Deflation is divided into three categories: (1) surface creep, where larger, heavier particles slide or roll along the ground; (2) saltation, where particles are lifted a short height into the air, and bounce and saltate across the surface of the soil; and (3) suspension, where very small and light particles are lifted into the air by the wind, and are often carried for long distances. Saltation is responsible for the majority (50–70%) of wind erosion, followed by suspension (30–40%), and then surface creep (5–25%). Silty soils tend to be the most affected by wind erosion; silt particles are relatively easily detached and carried away.
Wind erosion is much more severe in arid areas and during times of drought. For example, in the Great Plains, it is estimated that soil loss due to wind erosion can be as much as 6100 times greater in drought years than in wet years.
Mass movement
Mass movement is the downward and outward movement of rock and sediments on a sloped surface, mainly due to the force of gravity.
Mass movement is an important part of the erosional process, and is often the first stage in the breakdown and transport of weathered materials in mountainous areas. It moves material from higher elevations to lower elevations where other eroding agents such as streams and glaciers can then pick up the material and move it to even lower elevations. Mass-movement processes are always occurring continuously on all slopes; some mass-movement processes act very slowly; others occur very suddenly, often with disastrous results. Any perceptible down-slope movement of rock or sediment is often referred to in general terms as a landslide. However, landslides can be classified in a much more detailed way that reflects the mechanisms responsible for the movement and the velocity at which the movement occurs. One of the visible topographical manifestations of a very slow form of such activity is a scree slope.
Slumping happens on steep hillsides, occurring along distinct fracture zones, often within materials like clay that, once released, may move quite rapidly downhill. They will often show a spoon-shaped isostatic depression, in which the material has begun to slide downhill. In some cases, the slump is caused by water beneath the slope weakening it. In many cases it is simply the result of poor engineering along highways where it is a regular occurrence.
Surface creep is the slow movement of soil and rock debris by gravity which is usually not perceptible except through extended observation. However, the term can also describe the rolling of dislodged soil particles in diameter by wind along the soil surface.
Tillage erosion
Factors affecting soil erosion
Climate
The amount and intensity of precipitation is the main climatic factor governing soil erosion by water. The relationship is particularly strong if heavy rainfall occurs at times when, or in locations where, the soil's surface is not well protected by vegetation. This might be during periods when agricultural activities leave the soil bare, or in semi-arid regions where vegetation is naturally sparse. Wind erosion requires strong winds, particularly during times of drought when vegetation is sparse and soil is dry (and so is more erodible). Other climatic factors such as average temperature and temperature range may also affect erosion, via their effects on vegetation and soil properties. In general, given similar vegetation and ecosystems, areas with more precipitation (especially high-intensity rainfall), more wind, or more storms are expected to have more erosion.
In some areas of the world (e.g. the Midwestern United States and the Amazon Rainforest), rainfall intensity is the primary determinant of erosivity, with higher intensity rainfall generally resulting in more soil erosion by water. The size and velocity of rain drops is also an important factor. Larger and higher-velocity rain drops have greater kinetic energy, and thus their impact will displace soil particles by larger distances than smaller, slower-moving rain drops.
In other regions of the world (e.g. western Europe), runoff and erosion result from relatively low intensities of stratiform rainfall falling onto previously saturated soil. In such situations, rainfall amount rather than intensity is the main factor determining the severity of soil erosion by water.
Soil structure and composition
The composition, moisture, and compaction of soil are all major factors in determining the erosivity of rainfall. Sediments containing more clay tend to be more resistant to erosion than those with sand or silt, because the clay helps bind soil particles together. Soil containing high levels of organic materials are often more resistant to erosion, because the organic materials coagulate soil colloids and create a stronger, more stable soil structure. The amount of water present in the soil before the precipitation also plays an important role, because it sets limits on the amount of water that can be absorbed by the soil (and hence prevented from flowing on the surface as erosive runoff). Wet, saturated soils will not be able to absorb as much rainwater, leading to higher levels of surface runoff and thus higher erosivity for a given volume of rainfall. Soil compaction also affects the permeability of the soil to water, and hence the amount of water that flows away as runoff. More compacted soils will have a larger amount of surface runoff than less compacted soils.
Vegetative cover
Vegetation acts as an interface between the atmosphere and the soil. It increases the permeability of the soil to rainwater, thus decreasing runoff. It shelters the soil from winds, which results in decreased wind erosion, as well as advantageous changes in microclimate. The roots of the plants bind the soil together, and interweave with other roots, forming a more solid mass that is less susceptible to both water and wind erosion. The removal of vegetation increases the rate of surface erosion.
Topography
The topography of the land determines the velocity at which surface runoff will flow, which in turn determines the erosivity of the runoff. Longer, steeper slopes (especially those without adequate vegetative cover) are more susceptible to very high rates of erosion during heavy rains than shorter, less steep slopes. Steeper terrain is also more prone to mudslides, landslides, and other forms of gravitational erosion processes.
Human activities that aid soil erosion
Agricultural practices
Unsustainable agricultural practices increase rates of erosion by one to two orders of magnitude over the natural rate and far exceed replacement by soil production. The tillage of agricultural lands, which breaks up soil into finer particles, is one of the primary factors. The problem has been exacerbated in modern times, due to mechanized agricultural equipment that allows for deep plowing, which severely increases the amount of soil that is available for transport by water erosion. Others include monocropping, farming on steep slopes, pesticide and chemical fertilizer usage (which kill organisms that bind soil together), row-cropping, and the use of surface irrigation. A complex overall situation with respect to defining nutrient losses from soils, could arise as a result of the size selective nature of soil erosion events. Loss of total phosphorus, for instance, in the finer eroded fraction is greater relative to the whole soil. Extrapolating this evidence to predict subsequent behaviour within receiving aquatic systems, the reason is that this more easily transported material may support a lower solution P concentration compared to coarser sized fractions. Tillage also increases wind erosion rates, by dehydrating the soil and breaking it up into smaller particles that can be picked up by the wind. Exacerbating this is the fact that most of the trees are generally removed from agricultural fields, allowing winds to have long, open runs to travel over at higher speeds. Heavy grazing reduces vegetative cover and causes severe soil compaction, both of which increase erosion rates.
Deforestation
In an undisturbed forest, the mineral soil is protected by a layer of leaf litter and an humus that cover the forest floor. These two layers form a protective mat over the soil that absorbs the impact of rain drops. They are porous and highly permeable to rainfall, and allow rainwater to slow percolate into the soil below, instead of flowing over the surface as runoff. The roots of the trees and plants hold together soil particles, preventing them from being washed away. The vegetative cover acts to reduce the velocity of the raindrops that strike the foliage and stems before hitting the ground, reducing their kinetic energy. However it is the forest floor, more than the canopy, that prevents surface erosion. The terminal velocity of rain drops is reached in about . Because forest canopies are usually higher than this, rain drops can often regain terminal velocity even after striking the canopy. However, the intact forest floor, with its layers of leaf litter and organic matter, is still able to absorb the impact of the rainfall.
Deforestation causes increased erosion rates due to exposure of mineral soil by removing the humus and litter layers from the soil surface, removing the vegetative cover that binds soil together, and causing heavy soil compaction from logging equipment. Once trees have been removed by fire or logging, infiltration rates become high and erosion low to the degree the forest floor remains intact. Severe fires can lead to significant further erosion if followed by heavy rainfall.
Globally one of the largest contributors to erosive soil loss in the year 2006 is the slash and burn treatment of tropical forests. In a number of regions of the earth, entire sectors of a country have been rendered unproductive. For example, on the Madagascar high central plateau, comprising approximately ten percent of that country's land area, virtually the entire landscape is sterile of vegetation, with gully erosive furrows typically in excess of deep and wide. Shifting cultivation is a farming system which sometimes incorporates the slash and burn method in some regions of the world. This degrades the soil and causes the soil to become less and less fertile.
Roads and human impact
Human Impact has major effects on erosion processes—first by denuding the land of vegetative cover, altering drainage patterns, and compacting the soil during construction; and next by covering the land in an impermeable layer of asphalt or concrete that increases the amount of surface runoff and increases surface wind speeds. Much of the sediment carried in runoff from urban areas (especially roads) is highly contaminated with fuel, oil, and other chemicals. This increased runoff, in addition to eroding and degrading the land that it flows over, also causes major disruption to surrounding watersheds by altering the volume and rate of water that flows through them, and filling them with chemically polluted sedimentation. The increased flow of water through local waterways also causes a large increase in the rate of bank erosion.
Climate change
The warmer atmospheric temperatures observed over the past decades are expected to lead to a more vigorous hydrological cycle, including more extreme rainfall events. The rise in sea levels that has occurred as a result of climate change has also greatly increased coastal erosion rates.
Studies on soil erosion suggest that increased rainfall amounts and intensities will lead to greater rates of soil erosion. Thus, if rainfall amounts and intensities increase in many parts of the world as expected, erosion will also increase, unless amelioration measures are taken. Soil erosion rates are expected to change in response to changes in climate for a variety of reasons. The most direct is the change in the erosive power of rainfall. Other reasons include: a) changes in plant canopy caused by shifts in plant biomass production associated with moisture regime; b) changes in litter cover on the ground caused by changes in both plant residue decomposition rates driven by temperature and moisture dependent soil microbial activity as well as plant biomass production rates; c) changes in soil moisture due to shifting precipitation regimes and evapo-transpiration rates, which changes infiltration and runoff ratios; d) soil erodibility changes due to decrease in soil organic matter concentrations in soils that lead to a soil structure that is more susceptible to erosion and increased runoff due to increased soil surface sealing and crusting; e) a shift of winter precipitation from non-erosive snow to erosive rainfall due to increasing winter temperatures; f) melting of permafrost, which induces an erodible soil state from a previously non-erodible one; and g) shifts in land use made necessary to accommodate new climatic regimes.
Studies by Pruski and Nearing indicated that, other factors such as land use unconsidered, it is reasonable to expect approximately a 1.7% change in soil erosion for each 1% change in total precipitation under climate change. In recent studies, there are predicted increases of rainfall erosivity by 17% in the United States, by 18% in Europe, and globally 30 to 66%
Global environmental effects
Due to the severity of its ecological effects, and the scale on which it is occurring, erosion constitutes one of the most significant global environmental problems we face today.
Land degradation
Water and wind erosion are now the two primary causes of land degradation; combined, they are responsible for 84% of degraded acreage.
Each year, about 75 billion tons of soil is eroded from the land—a rate that is about 13–40 times as fast as the natural rate of erosion. Approximately 40% of the world's agricultural land is seriously degraded. According to the United Nations, an area of fertile soil the size of Ukraine is lost every year because of drought, deforestation and climate change. In Africa, if current trends of soil degradation continue, the continent might be able to feed just 25% of its population by 2025, according to UNU's Ghana-based Institute for Natural Resources in Africa.
Recent modeling developments have quantified rainfall erosivity at global scale using high temporal resolution (<30 min) and high fidelity rainfall recordings. The results is an extensive global data collection effort produced the Global Rainfall Erosivity Database (GloREDa) which includes rainfall erosivity for 3,625 stations and covers 63 countries. This first ever Global Rainfall Erosivity Database was used to develop a global erosivity map at 30 arc-seconds(~1 km) based on sophisticated geostatistical process. According to a new study published in Nature Communications, almost 36 billion tons of soil is lost every year due to water, and deforestation and other changes in land use make the problem worse. The study investigates global soil erosion dynamics by means of high-resolution spatially distributed modelling (c. 250 × 250 m cell size). The geo-statistical approach allows, for the first time, the thorough incorporation into a global soil erosion model of land use and changes in land use, the extent, types, spatial distribution of global croplands and the effects of different regional cropping systems.
The loss of soil fertility due to erosion is further problematic because the response is often to apply chemical fertilizers, which leads to further water and soil pollution, rather than to allow the land to regenerate.
Sedimentation of aquatic ecosystems
Soil erosion (especially from agricultural activity) is considered to be the leading global cause of diffuse water pollution, due to the effects of the excess sediments flowing into the world's waterways. The sediments themselves act as pollutants, as well as being carriers for other pollutants, such as attached pesticide molecules or heavy metals.
The effect of increased sediments loads on aquatic ecosystems can be catastrophic. Silt can smother the spawning beds of fish, by filling in the space between gravel on the stream bed. It also reduces their food supply, and causes major respiratory issues for them as sediment enters their gills. The biodiversity of aquatic plant and algal life is reduced, and invertebrates are also unable to survive and reproduce. While the sedimentation event itself might be relatively short-lived, the ecological disruption caused by the mass die off often persists long into the future.
One of the most serious and long-running water erosion problems worldwide is in the People's Republic of China, on the middle reaches of the Yellow River and the upper reaches of the Yangtze River. From the Yellow River, over 1.6 billion tons of sediment flows into the ocean each year. The sediment originates primarily from water erosion in the Loess Plateau region of the northwest.
Airborne dust pollution
Soil particles picked up during wind erosion of soil are a major source of air pollution, in the form of airborne particulates—"dust". These airborne soil particles are often contaminated with toxic chemicals such as pesticides or petroleum fuels, posing ecological and public health hazards when they later land, or are inhaled/ingested.
Dust from erosion acts to suppress rainfall and changes the sky color from blue to white, which leads to an increase in red sunsets. Dust events have been linked to a decline in the health of coral reefs across the Caribbean and Florida, primarily since the 1970s. Similar dust plumes originate in the Gobi desert, which combined with pollutants, spread large distances downwind, or eastward, into North America.
Monitoring, measuring and modelling soil erosion
Monitoring and modeling of erosion processes can help people better understand the causes of soil erosion, make predictions of erosion under a range of possible conditions, and plan the implementation of preventative and restorative strategies for erosion. However, the complexity of erosion processes and the number of scientific disciplines that must be considered to understand and model them (e.g. climatology, hydrology, geology, soil science, agriculture, chemistry, physics, etc.) makes accurate modelling challenging. Erosion models are also non-linear, which makes them difficult to work with numerically, and makes it difficult or impossible to scale up to making predictions about large areas from data collected by sampling smaller plots.
The most commonly used model for predicting soil loss from water erosion is the Universal Soil Loss Equation (USLE). This was developed in the 1960s and 1970s. It estimates the average annual soil loss A on a plot-sized area as:
A = RKLSCP
where R is the rainfall erosivity factor, K is the soil erodibility factor, L and S are topographic factors representing length and slope, C is the cover and management factor and P is the support practices factor.
Despite the USLE's plot-scale spatial basis, the model has often been used to estimate soil erosion on much larger areas, such as watersheds, continents, and globally. One major problem is that the USLE cannot simulate gully erosion, and so erosion from gullies is ignored in any USLE-based assessment of erosion. Yet erosion from gullies can be a substantial proportion (10–80%) of total erosion on cultivated and grazed land.
During the 50 years since the introduction of the USLE, many other soil erosion models have been developed. But because of the complexity of soil erosion and its constituent processes, all erosion models can only roughly approximate actual erosion rates when validated i.e. when model predictions are compared with real-world measurements of erosion. Thus new soil erosion models continue to be developed. Some of these remain USLE-based, e.g. the G2 model. Other soil erosion models have largely (e.g. the Water Erosion Prediction Project model) or wholly (e.g. RHEM, the Rangeland Hydrology and Erosion Model ) abandoned usage of USLE elements. Global studies continue to be based on the USLE.
On a smaller scale (e.g. for individual channels, dams, or spillways), there are erosion rate models available based on the critical shear stress of erosion as well as the erodibility of the soil. These can be measured using geotechnical engineering methods such as the hole erosion test or the jet erosion test.
Prevention and remediation
The most effective known method for erosion prevention is to increase vegetative cover on the land, which helps prevent both wind and water erosion. Terracing is an extremely effective means of erosion control, which has been practiced for thousands of years by people all over the world. Windbreaks (also called shelterbelts) are rows of trees and shrubs that are planted along the edges of agricultural fields, to shield the fields against winds. In addition to significantly reducing wind erosion, windbreaks provide many other benefits such as improved microclimates for crops (which are sheltered from the dehydrating and otherwise damaging effects of wind), habitat for beneficial bird species, carbon sequestration, and aesthetic improvements to the agricultural landscape. Traditional planting methods, such as mixed-cropping (instead of monocropping) and crop rotation, have also been shown to significantly reduce erosion rates. Crop residues play a role in the mitigation of erosion, because they reduce the impact of raindrops breaking up the soil particles. There is a higher potential for erosion when producing potatoes than when growing cereals, or oilseed crops. Forages have a fibrous root system, which helps combat erosion by anchoring the plants to the top layer of the soil, and covering the entirety of the field, as it is a non-row crop. In tropical coastal systems, properties of mangroves have been examined as a potential means to reduce soil erosion. Their complex root structures are known to help reduce wave damage from storms and flood impacts while binding and building soils. These roots can slow down water flow, leading to the deposition of sediments and reduced erosion rates. However, in order to maintain sediment balance, adequate mangrove forest width needs to be present.
| Physical sciences | Soil science | Earth science |
59438 | https://en.wikipedia.org/wiki/Thermal%20conductivity%20and%20resistivity | Thermal conductivity and resistivity | The thermal conductivity of a material is a measure of its ability to conduct heat. It is commonly denoted by , , or and is measured in W·m−1·K−1.
Heat transfer occurs at a lower rate in materials of low thermal conductivity than in materials of high thermal conductivity. For instance, metals typically have high thermal conductivity and are very efficient at conducting heat, while the opposite is true for insulating materials such as mineral wool or Styrofoam. Correspondingly, materials of high thermal conductivity are widely used in heat sink applications, and materials of low thermal conductivity are used as thermal insulation. The reciprocal of thermal conductivity is called thermal resistivity.
The defining equation for thermal conductivity is , where is the heat flux, is the thermal conductivity, and is the temperature gradient. This is known as Fourier's law for heat conduction. Although commonly expressed as a scalar, the most general form of thermal conductivity is a second-rank tensor. However, the tensorial description only becomes necessary in materials which are anisotropic.
Definition
Simple definition
Consider a solid material placed between two environments of different temperatures. Let be the temperature at and be the temperature at , and suppose . An example of this scenario is a building on a cold winter day; the solid material in this case is the building wall, separating the cold outdoor environment from the warm indoor environment.
According to the second law of thermodynamics, heat will flow from the hot environment to the cold one as the temperature difference is equalized by diffusion. This is quantified in terms of a heat flux , which gives the rate, per unit area, at which heat flows in a given direction (in this case minus x-direction). In many materials, is observed to be directly proportional to the temperature difference and inversely proportional to the separation distance :
The constant of proportionality is the thermal conductivity; it is a physical property of the material. In the present scenario, since heat flows in the minus x-direction and is negative, which in turn means that . In general, is always defined to be positive. The same definition of can also be extended to gases and liquids, provided other modes of energy transport, such as convection and radiation, are eliminated or accounted for.
The preceding derivation assumes that the does not change significantly as temperature is varied from to . Cases in which the temperature variation of is non-negligible must be addressed using the more general definition of discussed below.
General definition
Thermal conduction is defined as the transport of energy due to random molecular motion across a temperature gradient. It is distinguished from energy transport by convection and molecular work in that it does not involve macroscopic flows or work-performing internal stresses.
Energy flow due to thermal conduction is classified as heat and is quantified by the vector , which gives the heat flux at position and time . According to the second law of thermodynamics, heat flows from high to low temperature. Hence, it is reasonable to postulate that is proportional to the gradient of the temperature field , i.e.
where the constant of proportionality, , is the thermal conductivity. This is called Fourier's law of heat conduction. Despite its name, it is not a law but a definition of thermal conductivity in terms of the independent physical quantities and . As such, its usefulness depends on the ability to determine for a given material under given conditions. The constant itself usually depends on and thereby implicitly on space and time. An explicit space and time dependence could also occur if the material is inhomogeneous or changing with time.
In some solids, thermal conduction is anisotropic, i.e. the heat flux is not always parallel to the temperature gradient. To account for such behavior, a tensorial form of Fourier's law must be used:
where is symmetric, second-rank tensor called the thermal conductivity tensor.
An implicit assumption in the above description is the presence of local thermodynamic equilibrium, which allows one to define a temperature field . This assumption could be violated in systems that are unable to attain local equilibrium, as might happen in the presence of strong nonequilibrium driving or long-ranged interactions.
Other quantities
In engineering practice, it is common to work in terms of quantities which are derivative to thermal conductivity and implicitly take into account design-specific features such as component dimensions.
For instance, thermal conductance is defined as the quantity of heat that passes in unit time through a plate of particular area and thickness when its opposite faces differ in temperature by one kelvin. For a plate of thermal conductivity , area and thickness , the conductance is , measured in W⋅K−1. The relationship between thermal conductivity and conductance is analogous to the relationship between electrical conductivity and electrical conductance.
Thermal resistance is the inverse of thermal conductance. It is a convenient measure to use in multicomponent design since thermal resistances are additive when occurring in series.
There is also a measure known as the heat transfer coefficient: the quantity of heat that passes per unit time through a unit area of a plate of particular thickness when its opposite faces differ in temperature by one kelvin. In ASTM C168-15, this area-independent quantity is referred to as the "thermal conductance". The reciprocal of the heat transfer coefficient is thermal insulance. In summary, for a plate of thermal conductivity , area and thickness ,
thermal conductance = , measured in W⋅K−1.
thermal resistance = , measured in K⋅W−1.
heat transfer coefficient = , measured in W⋅K−1⋅m−2.
thermal insulance = , measured in K⋅m2⋅W−1.
The heat transfer coefficient is also known as thermal admittance in the sense that the material may be seen as admitting heat to flow.
An additional term, thermal transmittance, quantifies the thermal conductance of a structure along with heat transfer due to convection and radiation. It is measured in the same units as thermal conductance and is sometimes known as the composite thermal conductance. The term U-value is also used.
Finally, thermal diffusivity combines thermal conductivity with density and specific heat:
.
As such, it quantifies the thermal inertia of a material, i.e. the relative difficulty in heating a material to a given temperature using heat sources applied at the boundary.
Units
In the International System of Units (SI), thermal conductivity is measured in watts per meter-kelvin (W/(m⋅K)). Some papers report in watts per centimeter-kelvin [W/(cm⋅K)].
However, physicists use other convenient units as well, e.g., in cgs units, where esu/(cm-sec-K) is used.
The Lorentz number, defined as L=κ/σT is a quantity independent of the carrier density and the scattering mechanism. Its value for a gas of non-interacting electrons (typical carriers in good metallic conductors) is 2.72×10−13 esu/K2, or equivalently, 2.44×10−8 Watt-Ohm/K2.
In imperial units, thermal conductivity is measured in BTU/(h⋅ft⋅°F).
The dimension of thermal conductivity is M1L1T−3Θ−1, expressed in terms of the dimensions mass (M), length (L), time (T), and temperature (Θ).
Other units which are closely related to the thermal conductivity are in common use in the construction and textile industries. The construction industry makes use of measures such as the R-value (resistance) and the U-value (transmittance or conductance). Although related to the thermal conductivity of a material used in an insulation product or assembly, R- and U-values are measured per unit area, and depend on the specified thickness of the product or assembly.
Likewise the textile industry has several units including the tog and the clo which express thermal resistance of a material in a way analogous to the R-values used in the construction industry.
Measurement
There are several ways to measure thermal conductivity; each is suitable for a limited range of materials. Broadly speaking, there are two categories of measurement techniques: steady-state and transient. Steady-state techniques infer the thermal conductivity from measurements on the state of a material once a steady-state temperature profile has been reached, whereas transient techniques operate on the instantaneous state of a system during the approach to steady state. Lacking an explicit time component, steady-state techniques do not require complicated signal analysis (steady state implies constant signals). The disadvantage is that a well-engineered experimental setup is usually needed, and the time required to reach steady state precludes rapid measurement.
In comparison with solid materials, the thermal properties of fluids are more difficult to study experimentally. This is because in addition to thermal conduction, convective and radiative energy transport are usually present unless measures are taken to limit these processes. The formation of an insulating boundary layer can also result in an apparent reduction in the thermal conductivity.
Experimental values
The thermal conductivities of common substances span at least four orders of magnitude. Gases generally have low thermal conductivity, and pure metals have high thermal conductivity. For example, under standard conditions the thermal conductivity of copper is over times that of air.
Of all materials, allotropes of carbon, such as graphite and diamond, are usually credited with having the highest thermal conductivities at room temperature. The thermal conductivity of natural diamond at room temperature is several times higher than that of a highly conductive metal such as copper (although the precise value varies depending on the diamond type).
Thermal conductivities of selected substances are tabulated below; an expanded list can be found in the list of thermal conductivities. These values are illustrative estimates only, as they do not account for measurement uncertainties or variability in material definitions.
Influencing factors
Temperature
The effect of temperature on thermal conductivity is different for metals and nonmetals. In metals, heat conductivity is primarily due to free electrons. Following the Wiedemann–Franz law, thermal conductivity of metals is approximately proportional to the absolute temperature (in kelvins) times electrical conductivity. In pure metals the electrical conductivity decreases with increasing temperature and thus the product of the two, the thermal conductivity, stays approximately constant. However, as temperatures approach absolute zero, the thermal conductivity decreases sharply. In alloys the change in electrical conductivity is usually smaller and thus thermal conductivity increases with temperature, often proportionally to temperature. Many pure metals have a peak thermal conductivity between 2 K and 10 K.
On the other hand, heat conductivity in nonmetals is mainly due to lattice vibrations (phonons). Except for high-quality crystals at low temperatures, the phonon mean free path is not reduced significantly at higher temperatures. Thus, the thermal conductivity of nonmetals is approximately constant at high temperatures. At low temperatures well below the Debye temperature, thermal conductivity decreases, as does the heat capacity, due to carrier scattering from defects.
Chemical phase
When a material undergoes a phase change (e.g. from solid to liquid), the thermal conductivity may change abruptly. For instance, when ice melts to form liquid water at 0 °C, the thermal conductivity changes from 2.18 W/(m⋅K) to 0.56 W/(m⋅K).
Even more dramatically, the thermal conductivity of a fluid diverges in the vicinity of the vapor-liquid critical point.
Thermal anisotropy
Some substances, such as non-cubic crystals, can exhibit different thermal conductivities along different crystal axes. Sapphire is a notable example of variable thermal conductivity based on orientation and temperature, with 35 W/(m⋅K) along the c axis and 32 W/(m⋅K) along the a axis.
Wood generally conducts better along the grain than across it. Other examples of materials where the thermal conductivity varies with direction are metals that have undergone heavy cold pressing, laminated materials, cables, the materials used for the Space Shuttle thermal protection system, and fiber-reinforced composite structures.
When anisotropy is present, the direction of heat flow may differ from the direction of the thermal gradient.
Electrical conductivity
In metals, thermal conductivity is approximately correlated with electrical conductivity according to the Wiedemann–Franz law, as freely moving valence electrons transfer not only electric current but also heat energy. However, the general correlation between electrical and thermal conductance does not hold for other materials, due to the increased importance of phonon carriers for heat in non-metals. Highly electrically conductive silver is less thermally conductive than diamond, which is an electrical insulator but conducts heat via phonons due to its orderly array of atoms.
Magnetic field
The influence of magnetic fields on thermal conductivity is known as the thermal Hall effect or Righi–Leduc effect.
Gaseous phases
In the absence of convection, air and other gases are good insulators. Therefore, many insulating materials function simply by having a large number of gas-filled pockets which obstruct heat conduction pathways. Examples of these include expanded and extruded polystyrene (popularly referred to as "styrofoam") and silica aerogel, as well as warm clothes. Natural, biological insulators such as fur and feathers achieve similar effects by trapping air in pores, pockets, or voids.
Low density gases, such as hydrogen and helium typically have high thermal conductivity. Dense gases such as xenon and dichlorodifluoromethane have low thermal conductivity. An exception, sulfur hexafluoride, a dense gas, has a relatively high thermal conductivity due to its high heat capacity. Argon and krypton, gases denser than air, are often used in insulated glazing (double paned windows) to improve their insulation characteristics.
The thermal conductivity through bulk materials in porous or granular form is governed by the type of gas in the gaseous phase, and its pressure. At low pressures, the thermal conductivity of a gaseous phase is reduced, with this behaviour governed by the Knudsen number, defined as , where is the mean free path of gas molecules and is the typical gap size of the space filled by the gas. In a granular material corresponds to the characteristic size of the gaseous phase in the pores or intergranular spaces.
Isotopic purity
The thermal conductivity of a crystal can depend strongly on isotopic purity, assuming other lattice defects are negligible. A notable example is diamond: at a temperature of around 100 K the thermal conductivity increases from 10,000 W·m−1·K−1 for natural type IIa diamond (98.9% 12C), to 41,000 for 99.9% enriched synthetic diamond. A value of 200,000 is predicted for 99.999% 12C at 80 K, assuming an otherwise pure crystal. The thermal conductivity of 99% isotopically enriched cubic boron nitride is ~ 1400 W·m−1·K−1, which is 90% higher than that of natural boron nitride.
Molecular origins
The molecular mechanisms of thermal conduction vary among different materials, and in general depend on details of the microscopic structure and molecular interactions. As such, thermal conductivity is difficult to predict from first-principles. Any expressions for thermal conductivity which are exact and general, e.g. the Green-Kubo relations, are difficult to apply in practice, typically consisting of averages over multiparticle correlation functions. A notable exception is a monatomic dilute gas, for which a well-developed theory exists expressing thermal conductivity accurately and explicitly in terms of molecular parameters.
In a gas, thermal conduction is mediated by discrete molecular collisions. In a simplified picture of a solid, thermal conduction occurs by two mechanisms: 1) the migration of free electrons and 2) lattice vibrations (phonons). The first mechanism dominates in pure metals and the second in non-metallic solids. In liquids, by contrast, the precise microscopic mechanisms of thermal conduction are poorly understood.
Gases
In a simplified model of a dilute monatomic gas, molecules are modeled as rigid spheres which are in constant motion, colliding elastically with each other and with the walls of their container. Consider such a gas at temperature and with density , specific heat and molecular mass . Under these assumptions, an elementary calculation yields for the thermal conductivity
where is a numerical constant of order , is the Boltzmann constant, and is the mean free path, which measures the average distance a molecule travels between collisions. Since is inversely proportional to density, this equation predicts that thermal conductivity is independent of density for fixed temperature. The explanation is that increasing density increases the number of molecules which carry energy but decreases the average distance a molecule can travel before transferring its energy to a different molecule: these two effects cancel out. For most gases, this prediction agrees well with experiments at pressures up to about 10 atmospheres. At higher densities, the simplifying assumption that energy is only transported by the translational motion of particles no longer holds, and the theory must be modified to account for the transfer of energy across a finite distance at the moment of collision between particles, as well as the locally non-uniform density in a high density gas. This modification has been carried out, yielding Revised Enskog Theory, which predicts a density dependence of the thermal conductivity in dense gases.
Typically, experiments show a more rapid increase with temperature than (here, is independent of ). This failure of the elementary theory can be traced to the oversimplified "hard sphere" model, which both ignores the "softness" of real molecules, and the attractive forces present between real molecules, such as dispersion forces.
To incorporate more complex interparticle interactions, a systematic approach is necessary. One such approach is provided by Chapman–Enskog theory, which derives explicit expressions for thermal conductivity starting from the Boltzmann equation. The Boltzmann equation, in turn, provides a statistical description of a dilute gas for generic interparticle interactions. For a monatomic gas, expressions for derived in this way take the form
where is an effective particle diameter and is a function of temperature whose explicit form depends on the interparticle interaction law. For rigid elastic spheres, is independent of and very close to . More complex interaction laws introduce a weak temperature dependence. The precise nature of the dependence is not always easy to discern, however, as is defined as a multi-dimensional integral which may not be expressible in terms of elementary functions, but must be evaluated numerically. However, for particles interacting through a Mie potential (a generalisation of the Lennard-Jones potential) highly accurate correlations for in terms of reduced units have been developed.
An alternate, equivalent way to present the result is in terms of the gas viscosity , which can also be calculated in the Chapman–Enskog approach:
where is a numerical factor which in general depends on the molecular model. For smooth spherically symmetric molecules, however, is very close to , not deviating by more than for a variety of interparticle force laws. Since , , and are each well-defined physical quantities which can be measured independent of each other, this expression provides a convenient test of the theory. For monatomic gases, such as the noble gases, the agreement with experiment is fairly good.
For gases whose molecules are not spherically symmetric, the expression still holds. In contrast with spherically symmetric molecules, however, varies significantly depending on the particular form of the interparticle interactions: this is a result of the energy exchanges between the internal and translational degrees of freedom of the molecules. An explicit treatment of this effect is difficult in the Chapman–Enskog approach. Alternately, the approximate expression was suggested by Eucken, where is the heat capacity ratio of the gas.
The entirety of this section assumes the mean free path is small compared with macroscopic (system) dimensions. In extremely dilute gases this assumption fails, and thermal conduction is described instead by an apparent thermal conductivity which decreases with density. Ultimately, as the density goes to the system approaches a vacuum, and thermal conduction ceases entirely.
Liquids
The exact mechanisms of thermal conduction are poorly understood in liquids: there is no molecular picture which is both simple and accurate. An example of a simple but very rough theory is that of Bridgman, in which a liquid is ascribed a local molecular structure similar to that of a solid, i.e. with molecules located approximately on a lattice. Elementary calculations then lead to the expression
where is the Avogadro constant, is the volume of a mole of liquid, and is the speed of sound in the liquid. This is commonly called Bridgman's equation.
Metals
For metals at low temperatures the heat is carried mainly by the free electrons. In this case the mean velocity is the Fermi velocity which is temperature independent. The mean free path is determined by the impurities and the crystal imperfections which are temperature independent as well. So the only temperature-dependent quantity is the heat capacity c, which, in this case, is proportional to T. So
with k0 a constant. For pure metals, k0 is large, so the thermal conductivity is high. At higher temperatures the mean free path is limited by the phonons, so the thermal conductivity tends to decrease with temperature. In alloys the density of the impurities is very high, so l and, consequently k, are small. Therefore, alloys, such as stainless steel, can be used for thermal insulation.
Lattice waves, phonons, in dielectric solids
Heat transport in both amorphous and crystalline dielectric solids is by way of elastic vibrations of the lattice (i.e., phonons). This transport mechanism is theorized to be limited by the elastic scattering of acoustic phonons at lattice defects. This has been confirmed by the experiments of Chang and Jones on commercial glasses and glass ceramics, where the mean free paths were found to be limited by "internal boundary scattering" to length scales of 10−2 cm to 10−3 cm.
The phonon mean free path has been associated directly with the effective relaxation length for processes without directional correlation. If Vg is the group velocity of a phonon wave packet, then the relaxation length is defined as:
where t is the characteristic relaxation time. Since longitudinal waves have a much greater phase velocity than transverse waves, Vlong is much greater than Vtrans, and the relaxation length or mean free path of longitudinal phonons will be much greater. Thus, thermal conductivity will be largely determined by the speed of longitudinal phonons.
Regarding the dependence of wave velocity on wavelength or frequency (dispersion), low-frequency phonons of long wavelength will be limited in relaxation length by elastic Rayleigh scattering. This type of light scattering from small particles is proportional to the fourth power of the frequency. For higher frequencies, the power of the frequency will decrease until at highest frequencies scattering is almost frequency independent. Similar arguments were subsequently generalized to many glass forming substances using Brillouin scattering.
Phonons in the acoustical branch dominate the phonon heat conduction as they have greater energy dispersion and therefore a greater distribution of phonon velocities. Additional optical modes could also be caused by the presence of internal structure (i.e., charge or mass) at a lattice point; it is implied that the group velocity of these modes is low and therefore their contribution to the lattice thermal conductivity λL (L) is small.
Each phonon mode can be split into one longitudinal and two transverse polarization branches. By extrapolating the phenomenology of lattice points to the unit cells it is seen that the total number of degrees of freedom is 3pq when p is the number of primitive cells with q atoms/unit cell. From these only 3p are associated with the acoustic modes, the remaining 3p(q − 1) are accommodated through the optical branches. This implies that structures with larger p and q contain a greater number of optical modes and a reduced λL.
From these ideas, it can be concluded that increasing crystal complexity, which is described by a complexity factor CF (defined as the number of atoms/primitive unit cell), decreases λL. This was done by assuming that the relaxation time τ decreases with increasing number of atoms in the unit cell and then scaling the parameters of the expression for thermal conductivity in high temperatures accordingly.
Describing anharmonic effects is complicated because an exact treatment as in the harmonic case is not possible, and phonons are no longer exact eigensolutions to the equations of motion. Even if the state of motion of the crystal could be described with a plane wave at a particular time, its accuracy would deteriorate progressively with time. Time development would have to be described by introducing a spectrum of other phonons, which is known as the phonon decay. The two most important anharmonic effects are the thermal expansion and the phonon thermal conductivity.
Only when the phonon number ‹n› deviates from the equilibrium value ‹n›0, can a thermal current arise as stated in the following expression
where v is the energy transport velocity of phonons. Only two mechanisms exist that can cause time variation of ‹n› in a particular region. The number of phonons that diffuse into the region from neighboring regions differs from those that diffuse out, or phonons decay inside the same region into other phonons. A special form of the Boltzmann equation
states this. When steady state conditions are assumed the total time derivate of phonon number is zero, because the temperature is constant in time and therefore the phonon number stays also constant. Time variation due to phonon decay is described with a relaxation time (τ) approximation
which states that the more the phonon number deviates from its equilibrium value, the more its time variation increases. At steady state conditions and local thermal equilibrium are assumed we get the following equation
Using the relaxation time approximation for the Boltzmann equation and assuming steady-state conditions, the phonon thermal conductivity λL can be determined. The temperature dependence for λL originates from the variety of processes, whose significance for λL depends on the temperature range of interest. Mean free path is one factor that determines the temperature dependence for λL, as stated in the following equation
where Λ is the mean free path for phonon and denotes the heat capacity. This equation is a result of combining the four previous equations with each other and knowing that for cubic or isotropic systems and .
At low temperatures (< 10 K) the anharmonic interaction does not influence the mean free path and therefore, the thermal resistivity is determined only from processes for which q-conservation does not hold. These processes include the scattering of phonons by crystal defects, or the scattering from the surface of the crystal in case of high quality single crystal. Therefore, thermal conductance depends on the external dimensions of the crystal and the quality of the surface. Thus, temperature dependence of λL is determined by the specific heat and is therefore proportional to T3.
Phonon quasimomentum is defined as ℏq and differs from normal momentum because it is only defined within an arbitrary reciprocal lattice vector. At higher temperatures (10 K < T < Θ), the conservation of energy and quasimomentum , where q1 is wave vector of the incident phonon and q2, q3 are wave vectors of the resultant phonons, may also involve a reciprocal lattice vector G complicating the energy transport process. These processes can also reverse the direction of energy transport.
Therefore, these processes are also known as Umklapp (U) processes and can only occur when phonons with sufficiently large q-vectors are excited, because unless the sum of q2 and q3 points outside of the Brillouin zone the momentum is conserved and the process is normal scattering (N-process). The probability of a phonon to have energy E is given by the Boltzmann distribution . To U-process to occur the decaying phonon to have a wave vector q1 that is roughly half of the diameter of the Brillouin zone, because otherwise quasimomentum would not be conserved.
Therefore, these phonons have to possess energy of , which is a significant fraction of Debye energy that is needed to generate new phonons. The probability for this is proportional to , with . Temperature dependence of the mean free path has an exponential form . The presence of the reciprocal lattice wave vector implies a net phonon backscattering and a resistance to phonon and thermal transport resulting finite λL, as it means that momentum is not conserved. Only momentum non-conserving processes can cause thermal resistance.
At high temperatures (T > Θ), the mean free path and therefore λL has a temperature dependence T−1, to which one arrives from formula by making the following approximation and writing . This dependency is known as Eucken's law and originates from the temperature dependency of the probability for the U-process to occur.
Thermal conductivity is usually described by the Boltzmann equation with the relaxation time approximation in which phonon scattering is a limiting factor. Another approach is to use analytic models or molecular dynamics or Monte Carlo based methods to describe thermal conductivity in solids.
Short wavelength phonons are strongly scattered by impurity atoms if an alloyed phase is present, but mid and long wavelength phonons are less affected. Mid and long wavelength phonons carry significant fraction of heat, so to further reduce lattice thermal conductivity one has to introduce structures to scatter these phonons. This is achieved by introducing interface scattering mechanism, which requires structures whose characteristic length is longer than that of impurity atom. Some possible ways to realize these interfaces are nanocomposites and embedded nanoparticles or structures.
Prediction
Because thermal conductivity depends continuously on quantities like temperature and material composition, it cannot be fully characterized by a finite number of experimental measurements. Predictive formulas become necessary if experimental values are not available under the physical conditions of interest. This capability is important in thermophysical simulations, where quantities like temperature and pressure vary continuously with space and time, and may encompass extreme conditions inaccessible to direct measurement.
In fluids
For the simplest fluids, such as monatomic gases and their mixtures at low to moderate densities, ab initio quantum mechanical computations can accurately predict thermal conductivity in terms of fundamental atomic properties—that is, without reference to existing measurements of thermal conductivity or other transport properties. This method uses Chapman-Enskog theory or Revised Enskog Theory to evaluate the thermal conductivity, taking fundamental intermolecular potentials as input, which are computed ab initio from a quantum mechanical description.
For most fluids, such high-accuracy, first-principles computations are not feasible. Rather, theoretical or empirical expressions must be fit to existing thermal conductivity measurements. If such an expression is fit to high-fidelity data over a large range of temperatures
and pressures, then it is called a "reference correlation" for that material. Reference correlations have been published for many pure materials; examples are carbon dioxide, ammonia, and benzene. Many of these cover temperature and pressure ranges that encompass gas, liquid, and supercritical phases.
Thermophysical modeling software often relies on reference correlations for predicting thermal conductivity at user-specified temperature and pressure. These correlations may be proprietary. Examples are REFPROP (proprietary) and CoolProp (open-source).
Thermal conductivity can also be computed using the Green-Kubo relations, which express transport coefficients in terms of the statistics of molecular trajectories. The advantage of these expressions is that they are formally exact and valid for general systems. The disadvantage is that they require detailed knowledge of particle trajectories, available only in computationally expensive simulations such as molecular dynamics. An accurate model for interparticle interactions is also required, which may be difficult to obtain for complex molecules.
History
Jan Ingenhousz and the thermal conductivity of different metals
In a 1780 letter to Benjamin Franklin, Dutch-born British scientist Jan Ingenhousz relates an experiment which enabled him to rank seven different metals according to their thermal conductivities:
| Physical sciences | Thermodynamics | Physics |
59442 | https://en.wikipedia.org/wiki/Baryte | Baryte | Baryte, barite or barytes ( or ) is a mineral consisting of barium sulfate (BaSO4). Baryte is generally white or colorless, and is the main source of the element barium. The baryte group consists of baryte, celestine (strontium sulfate), anglesite (lead sulfate), and anhydrite (calcium sulfate). Baryte and celestine form a solid solution .
Names and history
The radiating form, sometimes referred to as Bologna Stone, attained some notoriety among alchemists for specimens found in the 17th century near Bologna by Vincenzo Casciarolo. These became phosphorescent upon being calcined.
Carl Scheele determined that baryte contained a new element in 1774, but could not isolate barium, only barium oxide. Johan Gottlieb Gahn also isolated barium oxide two years later in similar studies. Barium was first isolated by electrolysis of molten barium salts in 1808 by Sir Humphry Davy in England.
The American Petroleum Institute specification API 13/ISO 13500, which governs baryte for drilling purposes, does not refer to any specific mineral, but rather a material that meets that specification. In practice, however, this is usually the mineral baryte.
The term "primary barytes" refers to the first marketable product, which includes crude baryte (run of mine) and the products of simple beneficiation methods, such as washing, jigging, heavy media separation, tabling, and flotation. Most crude baryte requires some upgrading to minimum purity or density. Baryte that is used as an aggregate in a "heavy" cement is crushed and screened to a uniform size. Most baryte is ground to a small, uniform size before it is used as a filler or extender, an addition to industrial products, in the production of barium chemicals or as a weighting agent in petroleum well drilling mud.
Name
The name baryte is derived from the , 'heavy'. The American spelling is barite. The International Mineralogical Association initially adopted "barite" as the official spelling, but recommended adopting the older "baryte" spelling later. This move was controversial and was notably ignored by American mineralogists.
Other names have been used for baryte, including barytine, barytite, barytes, heavy spar, tiff, and blanc fixe.
Mineral associations and locations
Baryte occurs in many depositional environments, and is deposited through many processes including biogenic, hydrothermal, and evaporation, among others. Baryte commonly occurs in lead-zinc veins in limestones, in hot spring deposits, and with hematite ore. It is often associated with the minerals anglesite and celestine. It has also been identified in meteorites.
Baryte has been found at locations in Australia, Brazil, Nigeria, Canada, Chile, China, India, Pakistan, Germany, Greece, Guatemala, Iran, Ireland (where it was mined on Benbulben), Liberia, Mexico, Morocco, Peru, Romania (Baia Sprie), Turkey, South Africa (Barberton Mountain Land), Thailand, United Kingdom (Cornwall, Cumbria, Dartmoor/Devon, Derbyshire, Durham, Shropshire, Perthshire, Argyllshire, and Surrey) and in the US from Cheshire, Connecticut, De Kalb, New York, and Fort Wallace, New Mexico. It is mined in Arkansas, Connecticut, Virginia, North Carolina, Georgia, Tennessee, Kentucky, Nevada, and Missouri.
The global production of baryte in 2019 was estimated to be around 9.5 million metric tons, down from 9.8 million metric tons in 2012. The major barytes producers (in thousand tonnes, data for 2017) are as follows: China (3,600), India (1,600), Morocco (1,000), Mexico (400), United States (330), Iran (280), Turkey (250), Russia (210), Kazakhstan (160), Thailand (130) and Laos (120).
The main users of barytes in 2017 were (in million tonnes) US (2.35), China (1.60), Middle East (1.55), the European Union and Norway (0.60), Russia and CIS (0.5), South America (0.35), Africa (0.25), and Canada (0.20). 70% of barytes was destined for oil and gas well drilling muds. 15% for barium chemicals, 14% for filler applications in automotive, construction, and paint industries, and 1% other applications.
Natural baryte formed under hydrothermal conditions may be associated with quartz or silica. In hydrothermal vents, the baryte-silica mineralisation can also be accompanied by precious metals.
Information about the mineral resource base of baryte ores is presented in some scientific articles.
Uses
In oil and gas drilling
Worldwide, 69–77% of baryte is used as a weighting agent for drilling fluids in oil and gas exploration to suppress high formation pressures and prevent blowouts. As a well is drilled, the bit passes through various formations, each with different characteristics. The deeper the hole, the more baryte is needed as a percentage of the total mud mix. An additional benefit of baryte is that it is non-magnetic and thus does not interfere with magnetic measurements taken in the borehole, either during logging-while-drilling or in separate drill hole logging. Baryte used for drilling petroleum wells can be black, blue, brown or gray depending on the ore body. The baryte is finely ground so that at least 97% of the material, by weight, can pass through a 200-mesh (75 μm) screen, and no more than 30%, by weight, can be less than 6 μm diameter. The ground baryte also must be dense enough so that its specific gravity is 4.2 or greater, soft enough to not damage the bearings of a tricone drill bit, chemically inert, and containing no more than 250 milligrams per kilogram of soluble alkaline salts. In August 2010, the American Petroleum Institute published specifications to modify the 4.2 drilling grade standards for baryte to include 4.1 SG materials.
In oxygen and sulfur isotopic analysis
In the deep ocean, away from continental sources of sediment, pelagic baryte precipitates and forms a significant amount of the sediments. Since baryte has oxygen, systematics in the δ18O of these sediments have been used to help constrain paleotemperatures for oceanic crust.
The variations in sulfur isotopes (34S/32S) are being examined in evaporite minerals containing sulfur (e.g. baryte) and carbonate associated sulfates (CAS) to determine past seawater sulfur concentrations which can help identify specific depositional periods such as anoxic or oxic conditions. The use of sulfur isotope reconstruction is often paired with oxygen when a molecule contains both elements.
Geochronological dating
Dating the baryte in hydrothermal vents has been one of the major methods to determine their ages. Common methods to date hydrothermal baryte include radiometric dating and electron spin resonance dating.
Other uses
Baryte is used in added-value applications which include filler in paint and plastics, sound reduction in engine compartments, coat of automobile finishes for smoothness and corrosion resistance, friction products for automobiles and trucks, radiation shielding concrete, glass ceramics, and medical applications (for example, a barium meal before a contrast CT scan). Baryte is supplied in a variety of forms and the price depends on the amount of processing; filler applications commanding higher prices following intense physical processing by grinding and micronising, and there are further premiums for whiteness and brightness and color. It is also used to produce other barium chemicals, notably barium carbonate which is used for the manufacture of LED glass for television and computer screens (historically in cathode-ray tubes); and for dielectrics.
Historically, baryte was used for the production of barium hydroxide for sugar refining, and as a white pigment for textiles, paper, and paint.
Although baryte contains the toxic alkaline earth metal barium, it is not detrimental for human health, animals, plants and the environment because barium sulfate is extremely insoluble in water.
It is also sometimes used as a gemstone.
| Physical sciences | Minerals | Earth science |
59444 | https://en.wikipedia.org/wiki/Energy%20level | Energy level | A quantum mechanical system or particle that is bound—that is, confined spatially—can only take on certain discrete values of energy, called energy levels. This contrasts with classical particles, which can have any amount of energy. The term is commonly used for the energy levels of the electrons in atoms, ions, or molecules, which are bound by the electric field of the nucleus, but can also refer to energy levels of nuclei or vibrational or rotational energy levels in molecules. The energy spectrum of a system with such discrete energy levels is said to be quantized.
In chemistry and atomic physics, an electron shell, or principal energy level, may be thought of as the orbit of one or more electrons around an atom's nucleus. The closest shell to the nucleus is called the "1 shell" (also called "K shell"), followed by the "2 shell" (or "L shell"), then the "3 shell" (or "M shell"), and so on further and further from the nucleus. The shells correspond with the principal quantum numbers ( = 1, 2, 3, 4, ...) or are labeled alphabetically with letters used in the X-ray notation (K, L, M, N, ...).
Each shell can contain only a fixed number of electrons: The first shell can hold up to two electrons, the second shell can hold up to eight (2 + 6) electrons, the third shell can hold up to 18 (2 + 6 + 10) and so on. The general formula is that the nth shell can in principle hold up to 2n2 electrons. Since electrons are electrically attracted to the nucleus, an atom's electrons will generally occupy outer shells only if the more inner shells have already been completely filled by other electrons. However, this is not a strict requirement: atoms may have two or even three incomplete outer shells. (See Madelung rule for more details.) For an explanation of why electrons exist in these shells see electron configuration.
If the potential energy is set to zero at infinite distance from the atomic nucleus or molecule, the usual convention, then bound electron states have negative potential energy.
If an atom, ion, or molecule is at the lowest possible energy level, it and its electrons are said to be in the ground state. If it is at a higher energy level, it is said to be excited, or any electrons that have higher energy than the ground state are excited. An energy level is regarded as degenerate if there is more than one measurable quantum mechanical state associated with it.
Explanation
Quantized energy levels result from the wave behavior of particles, which gives a relationship between a particle's energy and its wavelength. For a confined particle such as an electron in an atom, the wave functions that have well defined energies have the form of a standing wave. States having well-defined energies are called stationary states because they are the states that do not change in time. Informally, these states correspond to a whole number of wavelengths of the wavefunction along a closed path (a path that ends where it started), such as a circular orbit around an atom, where the number of wavelengths gives the type of atomic orbital (0 for s-orbitals, 1 for p-orbitals and so on). Elementary examples that show mathematically how energy levels come about are the particle in a box and the quantum harmonic oscillator.
Any superposition (linear combination) of energy states is also a quantum state, but such states change with time and do not have well-defined energies. A measurement of the energy results in the collapse of the wavefunction, which results in a new state that consists of just a single energy state. Measurement of the possible energy levels of an object is called spectroscopy.
History
The first evidence of quantization in atoms was the observation of spectral lines in light from the sun in the early 1800s by Joseph von Fraunhofer and William Hyde Wollaston. The notion of energy levels was proposed in 1913 by Danish physicist Niels Bohr in the Bohr theory of the atom. The modern quantum mechanical theory giving an explanation of these energy levels in terms of the Schrödinger equation was advanced by Erwin Schrödinger and Werner Heisenberg in 1926.
Atoms
Intrinsic energy levels
In the formulas for energy of electrons at various levels given below in an atom, the zero point for energy is set when the electron in question has completely left the atom; i.e. when the electron's principal quantum number . When the electron is bound to the atom in any closer value of , the electron's energy is lower and is considered negative.
Orbital state energy level: atom/ion with nucleus + one electron
Assume there is one electron in a given atomic orbital in a hydrogen-like atom (ion). The energy of its state is mainly determined by the electrostatic interaction of the (negative) electron with the (positive) nucleus. The energy levels of an electron around a nucleus are given by:
(typically between 1 eV and 103 eV), where is the Rydberg constant, is the atomic number, is the principal quantum number, is the Planck constant, and is the speed of light. For hydrogen-like atoms (ions) only, the Rydberg levels depend only on the principal quantum number .
This equation is obtained from combining the Rydberg formula for any hydrogen-like element (shown below) with assuming that the principal quantum number above = in the Rydberg formula and (principal quantum number of the energy level the electron descends from, when emitting a photon). The Rydberg formula was derived from empirical spectroscopic emission data.
An equivalent formula can be derived quantum mechanically from the time-independent Schrödinger equation with a kinetic energy Hamiltonian operator using a wave function as an eigenfunction to obtain the energy levels as eigenvalues, but the Rydberg constant would be replaced by other fundamental physics constants.
Electron–electron interactions in atoms
If there is more than one electron around the atom, electron–electron interactions raise the energy level. These interactions are often neglected if the spatial overlap of the electron wavefunctions is low.
For multi-electron atoms, interactions between electrons cause the preceding equation to be no longer accurate as stated simply with as the atomic number. A simple (though not complete) way to understand this is as a shielding effect, where the outer electrons see an effective nucleus of reduced charge, since the inner electrons are bound tightly to the nucleus and partially cancel its charge. This leads to an approximate correction where is substituted with an effective nuclear charge symbolized as that depends strongly on the principal quantum number.
In such cases, the orbital types (determined by the azimuthal quantum number ) as well as their levels within the molecule affect and therefore also affect the various atomic electron energy levels. The Aufbau principle of filling an atom with electrons for an electron configuration takes these differing energy levels into account. For filling an atom with electrons in the ground state, the lowest energy levels are filled first and consistent with the Pauli exclusion principle, the Aufbau principle, and Hund's rule.
Fine structure splitting
Fine structure arises from relativistic kinetic energy corrections, spin–orbit coupling (an electrodynamic interaction between the electron's spin and motion and the nucleus's electric field) and the Darwin term (contact term interaction of shell electrons inside the nucleus). These affect the levels by a typical order of magnitude of 10−3 eV.
Hyperfine structure
This even finer structure is due to electron–nucleus spin–spin interaction, resulting in a typical change in the energy levels by a typical order of magnitude of 10−4 eV.
Energy levels due to external fields
Zeeman effect
There is an interaction energy associated with the magnetic dipole moment, , arising from the electronic orbital angular momentum, , given by
with
.
Additionally taking into account the magnetic momentum arising from the electron spin.
Due to relativistic effects (Dirac equation), there is a magnetic momentum, , arising from the electron spin
,
with the electron-spin g-factor (about 2), resulting in a total magnetic moment, ,
.
The interaction energy therefore becomes
.
Stark effect
Molecules
Chemical bonds between atoms in a molecule form because they make the situation more stable for the involved atoms, which generally means the sum energy level for the involved atoms in the molecule is lower than if the atoms were not so bonded. As separate atoms approach each other to covalently bond, their orbitals affect each other's energy levels to form bonding and antibonding molecular orbitals. The energy level of the bonding orbitals is lower, and the energy level of the antibonding orbitals is higher. For the bond in the molecule to be stable, the covalent bonding electrons occupy the lower energy bonding orbital, which may be signified by such symbols as σ or π depending on the situation. Corresponding anti-bonding orbitals can be signified by adding an asterisk to get σ* or π* orbitals. A non-bonding orbital in a molecule is an orbital with electrons in outer shells which do not participate in bonding and its energy level is the same as that of the constituent atom. Such orbitals can be designated as n orbitals. The electrons in an n orbital are typically lone pairs.
In polyatomic molecules, different vibrational and rotational energy levels are also involved.
Roughly speaking, a molecular energy state (i.e., an eigenstate of the molecular Hamiltonian) is the sum of the electronic, vibrational, rotational, nuclear, and translational components, such that:
where is an eigenvalue of the electronic molecular Hamiltonian (the value of the potential energy surface) at the equilibrium geometry of the molecule.
The molecular energy levels are labelled by the molecular term symbols. The specific energies of these components vary with the specific energy state and the substance.
Energy level diagrams
There are various types of energy level diagrams for bonds between atoms in a molecule.
Examples Molecular orbital diagrams, Jablonski diagrams, and Franck–Condon diagrams.
Energy level transitions
Electrons in atoms and molecules can change (make transitions in) energy levels by emitting or absorbing a photon (of electromagnetic radiation), whose energy must be exactly equal to the energy difference between the two levels.
Electrons can also be completely removed from a chemical species such as an atom, molecule, or ion. Complete removal of an electron from an atom can be a form of ionization, which is effectively moving the electron out to an orbital with an infinite principal quantum number, in effect so far away so as to have practically no more effect on the remaining atom (ion). For various types of atoms, there are 1st, 2nd, 3rd, etc. ionization energies for removing the 1st, then the 2nd, then the 3rd, etc. of the highest energy electrons, respectively, from the atom originally in the ground state. Energy in corresponding opposite quantities can also be released, sometimes in the form of photon energy, when electrons are added to positively charged ions or sometimes atoms. Molecules can also undergo transitions in their vibrational or rotational energy levels. Energy level transitions can also be nonradiative, meaning emission or absorption of a photon is not involved.
If an atom, ion, or molecule is at the lowest possible energy level, it and its electrons are said to be in the ground state. If it is at a higher energy level, it is said to be excited, or any electrons that have higher energy than the ground state are excited. Such a species can be excited to a higher energy level by absorbing a photon whose energy is equal to the energy difference between the levels. Conversely, an excited species can go to a lower energy level by spontaneously emitting a photon equal to the energy difference. A photon's energy is equal to the Planck constant () times its frequency () and thus is proportional to its frequency, or inversely to its wavelength ().
,
since , the speed of light, equals to
Correspondingly, many kinds of spectroscopy are based on detecting the frequency or wavelength of the emitted or absorbed photons to provide information on the material analyzed, including information on the energy levels and electronic structure of materials obtained by analyzing the spectrum.
An asterisk is commonly used to designate an excited state. An electron transition in a molecule's bond from a ground state to an excited state may have a designation such as σ → σ*, π → π*, or n → π* meaning excitation of an electron from a σ bonding to a σ antibonding orbital, from a π bonding to a π antibonding orbital, or from an n non-bonding to a π antibonding orbital. Reverse electron transitions for all these types of excited molecules are also possible to return to their ground states, which can be designated as σ* → σ, π* → π, or π* → n.
A transition in an energy level of an electron in a molecule may be combined with a vibrational transition and called a vibronic transition. A vibrational and rotational transition may be combined by rovibrational coupling. In rovibronic coupling, electron transitions are simultaneously combined with both vibrational and rotational transitions. Photons involved in transitions may have energy of various ranges in the electromagnetic spectrum, such as X-ray, ultraviolet, visible light, infrared, or microwave radiation, depending on the type of transition. In a very general way, energy level differences between electronic states are larger, differences between vibrational levels are intermediate, and differences between rotational levels are smaller, although there can be overlap. Translational energy levels are practically continuous and can be calculated as kinetic energy using classical mechanics.
Higher temperature causes fluid atoms and molecules to move faster increasing their translational energy, and thermally excites molecules to higher average amplitudes of vibrational and rotational modes (excites the molecules to higher internal energy levels). This means that as temperature rises, translational, vibrational, and rotational contributions to molecular heat capacity let molecules absorb heat and hold more internal energy. Conduction of heat typically occurs as molecules or atoms collide transferring the heat between each other. At even higher temperatures, electrons can be thermally excited to higher energy orbitals in atoms or molecules. A subsequent drop of an electron to a lower energy level can release a photon, causing a possibly coloured glow.
An electron further from the nucleus has higher potential energy than an electron closer to the nucleus, thus it becomes less bound to the nucleus, since its potential energy is negative and inversely dependent on its distance from the nucleus.
Crystalline materials
Crystalline solids are found to have energy bands, instead of or in addition to energy levels. Electrons can take on any energy within an unfilled band. At first this appears to be an exception to the requirement for energy levels. However, as shown in band theory, energy bands are actually made up of many discrete energy levels which are too close together to resolve. Within a band the number of levels is of the order of the number of atoms in the crystal, so although electrons are actually restricted to these energies, they appear to be able to take on a continuum of values. The important energy levels in a crystal are the top of the valence band, the bottom of the conduction band, the Fermi level, the vacuum level, and the energy levels of any defect states in the crystal.
| Physical sciences | Atomic physics | null |
59456 | https://en.wikipedia.org/wiki/Hydrate | Hydrate | In chemistry, a hydrate is a substance that contains water or its constituent elements. The chemical state of the water varies widely between different classes of hydrates, some of which were so labeled before their chemical structure was understood.
Chemical nature
Inorganic chemistry
Hydrates are inorganic salts "containing water molecules combined in a definite ratio as an integral part of the crystal" that are either bound to a metal center or that have crystallized with the metal complex. Such hydrates are also said to contain water of crystallization or water of hydration. If the water is heavy water in which the constituent hydrogen is the isotope deuterium, then the term deuterate may be used in place of hydrate.
A colorful example is cobalt(II) chloride, which turns from blue to red upon hydration, and can therefore be used as a water indicator.
The notation "hydrated compound⋅n", where n is the number of water molecules per formula unit of the salt, is commonly used to show that a salt is hydrated. The n is usually a low integer, though it is possible for fractional values to occur. For example, in a monohydrate n = 1, and in a hexahydrate n = 6. Numerical prefixes mostly of Greek origin are:
A hydrate that has lost water is referred to as an anhydride; the remaining water, if any exists, can only be removed with very strong heating. A substance that does not contain any water is referred to as anhydrous. Some anhydrous compounds are hydrated so easily that they are said to be hygroscopic and are used as drying agents or desiccants.
Organic chemistry
In organic chemistry, a hydrate is a compound formed by the hydration, i.e. "Addition of water or of the elements of water (i.e. H and OH) to a molecular entity". For example: ethanol, , is the product of the hydration reaction of ethene, , formed by the addition of H to one C and OH to the other C, and so can be considered as the hydrate of ethene. A molecule of water may be eliminated, for example, by the action of sulfuric acid. Another example is chloral hydrate, , which can be formed by reaction of water with chloral, .
Many organic molecules, as well as inorganic molecules, form crystals that incorporate water into the crystalline structure without chemical alteration of the organic molecule (water of crystallization). The sugar trehalose, for example, exists in both an anhydrous form (melting point 203 °C) and as a dihydrate (melting point 97 °C). Protein crystals commonly have as much as 50% water content.
Molecules are also labeled as hydrates for historical reasons not covered above. Glucose, , was originally thought of as and described as a carbohydrate.
Hydrate formation is common for active ingredients. Many manufacturing processes provide an opportunity for hydrates to form and the state of hydration can be changed with environmental humidity and time. The state of hydration of an active pharmaceutical ingredient can significantly affect the solubility and dissolution rate and therefore its bioavailability.
Clathrate hydrates
Clathrate hydrates (also known as gas hydrates, gas clathrates, etc.) are water ice with gas molecules trapped within; they are a form of clathrate. An important example is methane hydrate (also known as gas hydrate, methane clathrate, etc.).
Nonpolar molecules such as methane can form clathrate hydrates with water, especially under high pressure. Although there is no hydrogen bonding between water and guest molecules when methane is the guest molecule of the clathrate, guest–host hydrogen bonding often forms when the guest is a larger organic molecule such as tetrahydrofuran. In such cases the guest–host hydrogen bonds result in the formation of L-type Bjerrum defects in the clathrate lattice.
Stability
The stability of hydrates is generally determined by the nature of the compounds, their temperature, and the relative humidity (if they are exposed to air).
| Physical sciences | Water | Chemistry |
59493 | https://en.wikipedia.org/wiki/Fireworks | Fireworks | Fireworks are low explosive pyrotechnic devices used for aesthetic and entertainment purposes. They are most commonly used in fireworks displays (also called a fireworks show or pyrotechnics), combining a large number of devices in an outdoor setting. Such displays are the focal point of many cultural and religious celebrations, though mismanagement could lead to fireworks accidents.
Fireworks take many forms to produce four primary effects: noise, light, smoke, and floating materials (confetti most notably). They may be designed to burn with colored flames and sparks including red, orange, yellow, green, blue, purple and silver. They are generally classified by where they perform, either 'ground' or 'aerial'. Aerial fireworks may have their own propulsion (skyrocket) or be shot into the air by a mortar (aerial shell).
Most fireworks consist of a paper or pasteboard tube or casing filled with the combustible material, often pyrotechnic stars. A number of these tubes or cases may be combined so as to make when kindled, a great variety of sparkling shapes, often variously colored.
A skyrocket is a common form of firework, although the first skyrockets were used in warfare. The aerial shell, however, is the backbone of today's commercial aerial display, and a smaller version for consumer use is known as the festival ball in the United States.
Fireworks were originally invented in China. China remains the largest manufacturer and exporter of fireworks in the world.
Terminology
Silent fireworks
'Silent' fireworks displays are becoming popular due to concerns that noise effects traumatize pets, wildlife, and some humans. However, these are not a new type of firework and they are not completely silent. "Silent firework displays" refers to displays which simply exclude large, spectacular, noisy fireworks and make greater use of smaller, quieter devices.
History
The earliest fireworks came from China during the Song dynasty (960–1279). Fireworks were used to accompany many festivities. In China, pyrotechnicians were respected for their knowledge of complex techniques in creating fireworks and mounting firework displays.
During the Han dynasty (202 BC – 220 AD), people threw bamboo stems into a fire to produce an explosion with a loud sound. In later times, gunpowder packed into small containers was used to mimic the sounds of burning bamboo. Exploding bamboo stems and gunpowder firecrackers were interchangeably known as baozhu (爆竹) or baogan (爆竿). During the Song dynasty, people manufactured the first firecrackers comprising tubes made from rolled sheets of paper containing gunpowder and a fuse. They also strung these firecrackers together into large clusters, known as bian (lit. "whip") or bianpao (lit. "whip cannon"), so the firecrackers could be set off one by one in close sequence. By the 12th and possibly the 11th century, the term baozhang (爆仗) was used to specifically refer to gunpowder firecrackers. The first usage of the term was in the Dreams of the Glories of the Eastern Capital (東京夢華錄; about 1148) by Meng Yuanlao.
During the Song dynasty, common folk could purchase fireworks such as firecrackers from market vendors. Grand displays of fireworks were also known to be held. In 1110, according to the Dreams of the Glories of the Eastern Capital, a large fireworks display mounted by the military was held to entertain Emperor Huizong of Song (r. 1100–1125). The Qidong Yeyu (齊東野語; 1264) states that a rocket-propelled firework called a dilaoshu (地老鼠; lit. "earth rat") went off near the Empress Dowager Gong Sheng and startled her during a feast held in her honor by her son Emperor Lizong of Song (r. 1224–1264). This type of firework was one of the earliest examples of rocket propulsion. Around 1280, a Syrian named Hasan al-Rammah wrote of rockets, fireworks, and other incendiaries, using terms that suggested he derived his knowledge from Chinese sources, such as his references to fireworks as "Chinese flowers".
Colored fireworks were developed from earlier (possibly Han dynasty or soon thereafter) Chinese application of chemical substances to create colored smoke and fire. Such application appears in the Huolongjing (14th century) and Wubeizhi (preface of 1621, printed 1628), which describes recipes, several of which used low-nitrate gunpowder, to create military signal smokes with various colors. In the Wubei Huolongjing (武備火龍經; Ming, completed after 1628), two formulas appears for firework-like signals, the sanzhangju (三丈菊) and baizhanglian (百丈蓮), that produces silver sparkles in the smoke. In the Huoxilüe (火戲略; 1753) by Zhao Xuemin (趙學敏), there are several recipes with low-nitrate gunpowder and other chemical substances to tint flames and smoke. These included, for instance, arsenical sulphide for yellow, copper acetate (verdigris) for green, lead carbonate for lilac-white, and mercurous chloride (calomel) for white. The Chinese pyrotechnics were described by the French author Antoine Caillot (1818): "It is certain that the variety of colours which the Chinese have the secret of giving to flame is the greatest mystery of their fireworks." Similarly, the English geographer Sir John Barrow (ca. 1797) wrote "The diversity of colours indeed with which the Chinese have the secret of cloathing fire seems to be the chief merit of their pyrotechny."
Fireworks were produced in Europe by the 14th century, becoming popular by the 17th century. Lev Izmailov, ambassador of Peter the Great, once reported from China: "They make such fireworks that no one in Europe has ever seen." In 1758, the Jesuit missionary Pierre Nicolas le Chéron d'Incarville, living in Beijing, wrote about the methods and composition of Chinese fireworks to the Paris Academy of Sciences, which published the account five years later. Amédée-François Frézier published his revised work Traité des feux d'artice pour le spectacle (Treatise on Fireworks) in 1747 (originally 1706), covering the recreational and ceremonial uses of fireworks, rather than their military uses. Music for the Royal Fireworks was composed by George Frideric Handel in 1749 to celebrate the Peace treaty of Aix-la-Chapelle, which had been declared the previous year.
"Prior to the nineteenth century and the advent of modern chemistry they [fireworks] must have been relatively dull and unexciting." Bertholet in 1786 discovered that oxidations with potassium chlorate resulted in a violet emission. Subsequent developments revealed that oxidations with the chlorates of barium, strontium, copper, and sodium result in intense emission of bright colors. The isolation of metallic magnesium and aluminium marked another breakthrough as these metals burn with an intense silvery light.
Pyrotechnic compounds
Colors in fireworks are usually generated by pyrotechnic stars—usually just called stars—which produce intense light when ignited. Stars contain four basic types of ingredients.
A fuel
An oxidizer—a compound that combines with the fuel to produce intense heat
Color-producing salts (when the fuel itself is not the colorant)
A binder which holds the pellet together.
Some of the more common color-producing compounds are tabulated here. The color of a compound in a firework will be the same as its color in a flame test (shown at right). Not all compounds that produce a colored flame are appropriate for coloring fireworks, however. Ideal colorants will produce a pure, intense color when present in moderate concentration.
The color of sparks is limited to red/orange, yellow/gold and white/silver. This is explained by light emission from an incandescent solid particle in contrast to the element-specific emission from the vapor phase of a flame. Light emitted from a solid particle is defined by black-body radiation. Low boiling metals can form sparks with an intensively colored glowing shell surrounding the basic particle. This is caused by vapor phase combustion of the metal.
The brightest stars, often called Mag Stars, are fueled by aluminium. Magnesium is rarely used in the fireworks industry due to its lack of ability to form a protective oxide layer. Often an alloy of both metals called magnalium is used.
Many of the chemicals used in the manufacture of fireworks are non-toxic, while many more have some degree of toxicity, can cause skin sensitivity, or exist in dust form and are thereby inhalation hazards. Still others are poisons if directly ingested or inhaled.
Common elements in pyrotechnics
The following table lists the principal elements used in modern pyrotechnics. Some elements are used in their elemental form such as particles of titanium, aluminium, iron, zirconium, and magnesium. These elements burn in the presence of air (O2) or oxidants (perchlorate, chlorate). Most elements in pyrotechnics are in the form of salts.
Types
Aerial fireworks
Roman candle
A Roman candle is a long tube containing several large stars which fire at a regular interval. These are commonly arranged in fan shapes or crisscrossing shapes, at a closer proximity to the audience. Some larger Roman candles contain small shells (bombettes) rather than stars.
Mine
A mine (a.k.a. pot à feu) is a firework that expels stars and/or other garnitures into the sky. Shot from a mortar like a shell, a mine consists of a canister with the lift charge on the bottom with the effects placed on top. Mines can project small reports, serpents, and small shells, as well as just stars. Although mines up to diameter appear on occasion, they are usually in diameter.
Skyrocket
Cake
A cake is a cluster of individual tubes linked by fuse that fires a series of aerial effects. Tube diameters can range in size from , and a single cake can have more than 1,000 shots. The variety of effects within individual cakes is often such that they defy descriptive titles and are instead given cryptic names such as "Bermuda Triangle", "Pyro Glyphics", "Waco Wakeup", and "Poisonous Spider", to name a few. Others are simply quantities of shells fused together in single-shot tubes.
Ground fireworks
Sparkler
Firecracker
Bang snap
List of public display effects
Crossette
A shell containing several large stars that travel a short distance before breaking apart into smaller stars, creating a crisscrossing grid-like effect. Strictly speaking, a crossette star should split into four pieces which fly off symmetrically, making a cross. Once limited to silver or gold effects, colored crossettes such as red, green, or white are now very common.
Chrysanthemum
A spherical break of colored stars, similar to a peony, but with stars that leave a visible trail of sparks.
Dahlia
Essentially the same as a peony shell, but with fewer and larger stars. These stars travel a longer-than-usual distance from the shell break before burning out. For instance, if a peony shell is made with a star size designed for a shell, it is then considered a dahlia. Some dahlia shells are cylindrical rather than spherical to allow for larger stars.
Diadem
A type of chrysanthemum or peony, with a center cluster of non-moving stars, normally of a contrasting color or effect.
Fish
Inserts that propel themselves rapidly away from the shell burst, often resembling fish swimming away.
Horsetail
Named for the shape of its break, this shell features heavy long-burning tailed stars that only travel a short distance from the shell burst before free-falling to the ground. Also known as a waterfall shell. Sometimes there is a glittering through the "waterfall".
Kamuro
Kamuro is a Japanese word meaning "boys haircut", which is what this shell resembles when fully exploded in the air. It is a dense burst of glittering silver or gold stars which leave a heavy glitter trail and shine bright in the night's sky.
Multi-break shells
A large shell containing several smaller shells of various sizes and types. The initial burst scatters the shells across the sky before they explode. Also called a bouquet shell. When a shell contains smaller shells of the same size and type, the effect is usually referred to as "Thousands". Very large bouquet shells (up to ) are frequently used in Japan.
Palm
A shell containing a relatively few large comet stars arranged in such a way as to burst with large arms or tendrils, producing a palm tree-like effect. Proper palm shells feature a thick rising tail that displays as the shell ascends, thereby simulating the tree trunk to further enhance the "palm tree" effect. One might also see a burst of color inside the palm burst (given by a small insert shell) to simulate coconuts.
Peony
A spherical break of colored stars that burn without a tail effect. The peony is the most commonly seen shell type.
Ring
A shell with stars specially arranged so as to create a ring. Variations include smiley faces, hearts, and clovers.
Salute
A shell intended to produce a loud report rather than a visual effect. Salute shells usually contain flash powder, producing a quick flash followed by a very loud report resembling military artillery. Titanium may be added to the flash powder mix to produce a cloud of bright sparks around the flash. Salutes are commonly used in large quantities during finales to create intense noise and brightness. They are often cylindrical in shape to allow for a larger payload of flash powder, but ball shapes are common and cheaper as well. Salutes are also called Maroons.
Spider
A shell containing a fast burning tailed or charcoal star that is burst very hard so that the stars travel in a straight and flat trajectory before slightly falling and burning out. This appears in the sky as a series of radial lines much like the legs of a spider.
Time Rain
An effect created by large, slow-burning stars within a shell that leave a trail of large glittering sparks behind and make a sizzling noise. The "time" refers to the fact that these stars burn away gradually, as opposed to the standard brocade "rain" effect where a large amount of glitter material is released at once.
Willow
A willow is similar to a chrysanthemum, but with long-burning silver or gold stars that produce a soft, dome-shaped weeping willow-like effect.
Farfalle
Farfalle is an effect in Italian fireworks with spinning silver sprays in the air.
Tourbillion
A tourbillion is similar to a farfalle but has spinning stars.
Audio effects
Bang
The bang is the most common effect in fireworks and sounds like artillery cannon being fired; technically a "report". Silent fireworks have all of the visual effects, however. The "salute" effect is even more pronounced and sometimes is banned.
Crackle
The firework produces a crackling sound.
Hummer
Tiny tube fireworks that are ejected into the air spinning with such force that they shred their outer coating, in doing so they whizz and hum.
Whistle
High pitched often very loud screaming and screeching created by the resonance of gas. This is caused by a very fast strobing (on/off burning stage) of the fuel. The rapid bursts of gas from the fuel vibrate the air many hundreds of times per second causing the familiar whistling sound. It is not, as is commonly thought, made in the conventional way that musical instruments are using specific tube shapes or apertures. Common whistle fuels contain benzoate or salicylate compounds and a suitable oxidizer such as potassium perchlorate.
Safety and environmental impact
Improper use of fireworks is dangerous, both to the person operating them (risks of burns and wounds) and to bystanders; in addition, they may start fires on landing. To prevent fireworks accidents, the use of fireworks is legally restricted in many countries. In such countries, display fireworks are restricted for use by professionals; smaller consumer versions may or may not be available to the public.
Effects on animals
Birds and animals, both domestic and wild, can be frightened by their noise, leading to them running away, often into danger, or hurting themselves on fences or in other ways in an attempt to escape the perceived danger.
Majority of dogs experience distress, fear and anxiety during fireworks. In 2016, following a petition signed by more than 100,000 Brits, House of Commons of the United Kingdom debated a motion to restrict firework use.
Fireworks also affect birds, especially larger birds like geese, eagles and others. According to a study by Max Planck Institute and Netherlands Institute of Ecology, many birds abruptly leave their sleeping sites on New Year's Eve, and some fly up to 500 km non-stop to get away from human settlements. On average, about 1000 times more birds are in flight on New Year's Eve than on other nights. Frightened birds also may abandon nests and not return to complete rearing their young. A scientific study from 2022 indicates that fireworks might have some sort of lasting effect on birds, with many birds spending more time to find food in the weeks after New Year's Eve fireworks.
Pollution
Fireworks produce smoke and dust that may contain residues of heavy metals, sulfur-coal compounds and some low concentration toxic chemicals. These by-products of fireworks combustion will vary depending on the mix of ingredients of a particular firework. (The color green, for instance, may be produced by adding the various compounds and salts of barium, some of which are toxic, and some of which are not.) Some fishers have noticed and reported to environmental authorities that firework residues can hurt fish and other water-life because some may contain toxic compounds (such as antimony sulfide or arsenic). This is a subject of much debate due to the fact that large-scale pollution from other sources makes it difficult to measure the amount of pollution that comes specifically from fireworks. The possible toxicity of any fallout may also be affected by the amount of black powder used, type of oxidizer, colors produced and launch method.
Perchlorate salts, when in solid form, dissolve and move rapidly in groundwater and surface water. Even in low concentrations in drinking water supplies, perchlorate ions are known to inhibit the uptake of iodine by the thyroid gland. As of 2010, there are no federal drinking water standards for perchlorates in the United States, but the
US Environmental Protection Agency has studied the impacts of perchlorates on the environment as well as drinking water.
Several U.S. states have enacted drinking water standard for perchlorates, including Massachusetts in 2006. California's legislature enacted AB 826, the Perchlorate Contamination Prevention Act of 2003, requiring California's Department of Toxic Substance Control (DTSC) to adopt regulations specifying best management practices for perchlorate-containing substances. The Perchlorate Best Management Practices were adopted on 31 December 2005 and became operative on 1 July 2006. California issued drinking water standards in 2007. Several other states, including Arizona, Maryland, Nevada, New Mexico, New York, and Texas have established non-enforceable, advisory levels for perchlorates.
The courts have also taken action with regard to perchlorate contamination. For example, in 2003, a federal district court in California found that Comprehensive Environmental Response, Compensation and Liability Act (CERCLA) applied because perchlorate is ignitable and therefore a "characteristic" hazardous waste.
Pollutants from fireworks raise concerns because of potential health risks associated with the products of combustion during the liquid phase and the solid phase after they have cooled as well as the gases produced, particularly the carbon monoxide and carbon dioxide. For persons with asthma or other respiratory conditions, the smoke from fireworks may aggravate existing health problems.
Pollution is also a concern because fireworks often contain heavy metals as source of color. However, gunpowder smoke and the solid residues are basic, and as such the cumulative effect of fireworks on acid rain is uncertain. What is not disputed is that most consumer fireworks leave behind a considerable amount of solid debris, including both readily biodegradable components as well as nondegradable plastic items. Concerns over pollution, consumer safety, and debris have restricted the sale and use of consumer fireworks in many countries. Professional displays, on the other hand, remain popular around the world.
Others argue that alleged concern over pollution from fireworks constitutes a red herring, since the amount of contamination from fireworks is minuscule in comparison to emissions from sources such as the burning of fossil fuels. In the US, some states and local governments restrict the use of fireworks in accordance with the Clean Air Act which allows laws relating to the prevention and control of outdoor air pollution to be enacted.
Some companies within the U.S. fireworks industry claim they are working with Chinese manufacturers to reduce and ultimately hope to eliminate of the pollutant perchlorate.
Government regulations around the world
Australia
Fireworks are illegal in most Australian states and territories, unless part of a display by a licensed pyrotechnician and with a permit. However Tasmania, ACT and Northern Territory allow consumer use with a permit (dependent on calendar date and circumstances). On 1 July for Territory Day one can freely use fireworks without a permit in the Northern Territory.
Small novelties such as party poppers and sparklers are legal for consumers across Australia.
On 24 August 2009, the ACT Government announced a complete ban on backyard fireworks.
Canada
The use, storage and sale of commercial-grade fireworks in Canada is licensed by Natural Resources Canada's Explosive Regulatory Division (ERD). Unlike their consumer counterpart, commercial-grade fireworks function differently, and come in a wide range of sizes from up to or more in diameter. Commercial grade fireworks require a Fireworks Operator Certificate (FOC), obtained from the ERD by completing a one-day safety course. There are two categories of FOC: one for pyrotechnics (those used on stage and in movies) and another for display fireworks (those used in dedicated fireworks shows). Each requires completion of its own course, although there are special categories of FOC which allow visiting operators to run their shows with the assistance of a Canadian supervisor.
The display fireworks FOC has two levels: assistant, and fully licensed. A fully licensed display fireworks operator can also be further endorsed for marine launch, flying saucers, and other more technically demanding fireworks displays.
The pyrotechnician FOC has three levels: pyrotechnician (which allows work under a supervisor), supervising pyrotechnician, and special effects pyrotechnician (which allows the fabrication of certain types of pyrotechnic devices). Additionally, a special effects pyrotechnician can be endorsed for the use of detonating cord.
Since commercial-grade fireworks are shells which are loaded into separate mortars by hand, there is danger in every stage of the setup. Setup of these fireworks involves the placement and securing of mortars on wooden or wire racks; loading of the shells; and if electronically firing, wiring and testing. The mortars are generally made of FRE (fiber-reinforced epoxy) or HDPE (high-density polyethylene). Older mortars made of sheet steel have been banned by most countries due to the problem of shrapnel produced during a misfire.
Setup of mortars in Canada for an oblong firing site require that a mortar be configured at an angle of 10 to 15 degrees down-range with a safety distance of at least down-range and surrounding the mortars, plus distance adjustments for wind speed and direction. In June 2007, the ERD approved circular firing sites for use with vertically fired mortars with a safety distance of at least radius, plus distance adjustments for wind speed and direction.
Loading of shells is a delicate process, and must be done with caution, and a loader must ensure not only the mortar is clean, but also make sure that no part of their body is directly over the mortar in case of a premature fire. Wiring the shells is a painstaking process; whether the shells are being fired manually or electronically, any "chain fusing" or wiring of electrical ignitors, care must be taken to prevent the fuse (an electrical match, often incorrectly called a squib) from igniting. If the setup is wired electrically, the electrical matches are usually plugged into a "firing rail" or "breakout box" that runs back to the main firing board; from there, the Firing Board is simply hooked up to a car battery, and can proceed with firing the show when ready.
Since commercial-grade fireworks are so much larger and more powerful, setup, and firing crews are always under great pressure to ensure they safely set up, fire, and clean up after a show.
Chile
In Chile, the manufacture, importation, possession and use of fireworks is prohibited to unauthorized individuals; only certified firework companies can legally use fireworks. As they are considered a type of explosive, offenders can in principle be tried before military courts, although this is unusual in practice.
China
European Union
The European Union's policy is aimed at harmonising and standardising the EU member states' policies on the regulation of production, transportation, sale, consumption and overall safety of fireworks across Europe.
Belgium
In Belgium, each municipality can decide how to regulate fireworks. During New Year's Eve, lighting fireworks without a licence is allowed in 35% of the 308 Flemish municipalities, in around 50% a permit from the burgemeester (mayor) is required, and around 14% of municipalities have banned consumer fireworks altogether.
Finland
In Finland those under 18 years old haven't been allowed to buy any fireworks since 2009. Safety goggles are required. The use of fireworks is generally allowed on the evening and night of New Year's Eve, 31 December. In some municipalities of Western Finland it is allowed to use fireworks without a fire station's permission on the last weekend of August. With the fire station's permission, fireworks can be used year-round.
Germany
In Germany, amateurs over 18 years old are allowed to buy and ignite fireworks of Category F2 for several hours on 31 December and 1 January; each German municipality is authorised to limit the number of hours this may last locally. The sale of Category F3 and F4 fireworks to consumers is prohibited. Lighting fireworks is forbidden near churches, hospitals, retirement homes and wooden or thatch-roofed buildings. All major German cities organise professional fireworks shows.
In addition to the previously existing regulations, there was a nationwide ban on the sale of category F2 fireworks to consumers on New Year's Eve 2020/2021 during the COVID-19 pandemic, with the aim to relieve the burden on hospitals by reducing the number of emergencies due to injuries caused by fireworks on New Year's Eve. On the 2024–2025 New Year's eve and day, five were killed and 100+ were injured due to unusually intense widespread criminality and negligence.
Italy
In 2015, the Italian town of Collecchio mandated silent fireworks, being among the first to make the switch without losing the beauty of the visual displays.
Netherlands
In the Netherlands, fireworks cannot be sold to anyone under the age of 16. It may only be sold during a period of three days before a new year. If one of these days is a Sunday, that day is excluded from sale and sale may commence one day earlier.
Republic of Ireland
In the Republic of Ireland, fireworks are illegal and possession is punishable by huge fines and/or prison. However, around Halloween a large amount of fireworks are set off, due to the ease of being able to purchase from Northern Ireland.
Sweden
In Sweden, fireworks can only be purchased and used by people 18 or older. Firecrackers used to be banned, but are now allowed under European Union fireworks policy.
Iceland
In Iceland, the Icelandic law states that anyone may purchase and use fireworks during a certain period around New Year's Eve. Most places that sell fireworks in Iceland make their own rules about age of buyers, usually it is around 16. The people of Reykjavík spend enormous sums of money on fireworks, most of which are fired as midnight approaches on 31 December. As a result, every New Year's Eve the city is lit up with fireworks displays.
New Zealand
Fireworks in New Zealand are available from 2 to 5 November, around Guy Fawkes Day, and may be purchased only by those 18 years of age and older (up from 14 years pre-2007). Despite the restriction on when fireworks may be sold, there is no restriction regarding when fireworks may be used. The types of fireworks available to the public are multi-shot "cakes", Roman candles, single shot shooters, ground and wall spinners, fountains, cones, sparklers, and various novelties, such as smoke bombs and Pharaoh's serpents. Consumer fireworks are also not allowed to be louder than 90 decibels.
Norway
In Norway, fireworks can only be purchased and used by people 18 or older. Sale is restricted to a few days before New Year's Eve. Rockets are not allowed.
United Kingdom
Fireworks in the UK have become more strictly regulated since 1997. Since 2005, the law has been harmonised gradually, in accordance with other EU member state laws.
Fireworks are mostly used in England, Scotland and Wales around Diwali (late October or early November), on Guy Fawkes Night, 5 November and on New Year's Eve. In the UK, responsibility for the safety of firework displays is shared between the Health and Safety Executive, fire brigades and local authorities. Currently, there is no national system of licensing for fireworks operators, but in order to purchase large display fireworks, operators must have licensed explosives storage and public liability insurance.
Fireworks cannot be sold to people under the age of 18 and are not permitted to be set off between 11pm and 7am with exceptions only for:
Bonfire Night (5 November) (permitted until midnight)
The Chinese New Year (permitted until 1am)
Diwali (permitted until 1am)
New Year (permitted on New Year's Eve until 1am on New Year's Day)
The maximum legal NEC (net explosive content) of a UK firework available to the public is two kilograms. Jumping jacks, strings of firecrackers, shell firing tubes, bangers and mini-rockets were all banned during the late 1990s. In 2004, single-shot air bombs and bottle rockets were banned, and rocket sizes were limited. From March 2008 any firework with more than 5% flashpowder per tube has been classified 1.3G. The aim of these measures was to eliminate "pocket money" fireworks, and to limit the disruptive effects of loud bangs.
United States
In the United States, fireworks laws vary widely from state to state or county to county. Federal, state, and local authorities govern the use of display fireworks in the United States. The Consumer Product Safety Commission (CPSC) regulates consumer fireworks at the federal level through the Federal Hazardous Substances Act (FHSA). The National Fire Protection Association (NFPA) sets forth a set of codes that give the minimum standards of display fireworks use and safety in the U.S. Both state and local jurisdictions can further add restrictions on the use and safety requirements of display fireworks. There are currently 46 states in the United States in which fireworks are legal for consumer use.
Fireworks celebrations throughout the world
Australia
In Australia, fireworks displays are frequently used in the celebration of public holidays, particularly New Year's Eve and Australia Day. The most famous is the Sydney New Year's Eve Midnight Fireworks. In the Northern Territory, "Cracker Night" is celebrated every 1 July on Territory Day, where residents are allowed to buy and use fireworks without a permit.
Catalonia
Fireworks are an essential element of popular festival in Catalonia, especially the patron saint day for each Catalan town or city, usually called Festa Major. Coet (rocket) is the generic term for all kinds of pyrotechnic devices. Professional aerial displays are less common than the use of ground based fireworks by members of ritual crews.
The Correfoc (firerun) is an element of many festivals, in which crews of diables (devils) dance through the streets to beating drums, holding maces above their heads upon which are mounted carretilles which spin, producing a shower of sparks culminating in an explosive pop. There are three types, les normals, les xiuladores, which emit a whistling sound, and les Arboç, which produce a large umbrella of sparks. Those that emit sparks but do not spin are called francesos and a larger, more potent version are called portuguessos. Those that emit only light are simulators and those that serve to ignite all types of coets are called botafoc, only used by the cap de la colla, or crew chief, who decides when ignition occurs. In many cases, the devils light their carretilles simultaneously by holding all of their maces together and once ignited, begin dancing to the beat of pounding drums. The character of Llucifer carries a larger and more elaborate mace called a ceptrot which is said to "dominate over all," pictured at right.
Pyrotechnics feature in several other Catalan festes:
Nit de San Joan (St. John’s Night, June 14) townsfolk set off rockets from all over town and light bonfires.
Parliament de Diables (the infernal parliament) is a ritual performance by crews of devils in which characters such as Lucifer, the Diablessa (portrayed by a man in women’s clothes), and other members of the crew declaim satiric verses commenting on current events, punctuated by pyrotechnic dances.
La Patum is celebrated in Berga on Corpus Christi, which has featured a ball de diables called Els Plens since 1628. In 2005, UNESCO declared La Patum one of the Masterpieces of the Oral and Intangible Heritage of Humanity.
The correfoc of the festa major of Vilanova i la Geltrú culminates in a sostre de foc (ceiling of fire) during which a mass of coets attached to wires are strung above the central plaza enabling the assembled multitude to be showered with sparks and explosions from the sky.
France
In France, fireworks are traditionally displayed on the eve of Bastille day (14 July) to commemorate the French revolution and the storming of the Bastille on that same day in 1789. Every city in France lights up the sky for the occasion with a special mention to Paris that offers a spectacle around the Eiffel Tower.
Hungary
In Hungary fireworks are used on 20 August, which is a national celebration day
India
Indians throughout the world celebrate with fireworks as part of their popular "festival of lights" (Diwali) in Oct-Nov every year.
Japan
During the summer in Japan, are held nearly every day someplace in the country, numbering more than 200 during August alone. The festivals consist of large fireworks shows, the largest of which use between 100,000 and 120,000 rounds (Tondabayashi, Osaka), and can attract more than 800,000 spectators. Street vendors set up stalls to sell various drinks and staple Japanese food (such as yakisoba, okonomiyaki, takoyaki, kakigōri (shaved ice), and traditionally held festival games, such as kingyo-sukui, or goldfish scooping.
Even today, men and women attend these events wearing the traditional yukata, summer kimono, or jinbei, and gather in large social circles of family or friends to sit picnic-like, eating and drinking, while watching the show.
The first fireworks festival in Japan was held in 1733 when Tokugawa Yoshimune, the eight shogun of the Edo Period, decreed a festival be held on the Sumida River on July 9th every year to mourn those who died in famine and other disasters. This event is recognized as the start of the Sumidagawa Fireworks Festival, and is one of the most recognizable fireworks festivals in Japan.
Malta
Fireworks have been used in Malta for hundreds of years. When the islands were ruled by the Order of St John, fireworks were used on special occasions such as the election of a new Grand Master, the appointment of a new Pope or the birth of a prince.
Nowadays, fireworks are used in village feasts throughout the summer. The Malta International Fireworks Festival is also held annually.
Monte-Carlo International Fireworks Festival
Pyrotechnics experts from around the world have competed in Monte Carlo, Monaco, since 1966. The festival runs from July to August every year, and the winner returns in 18 November for the fireworks display on the night before the National Day of Monaco. The event is held in Port Hercule, beginning at around 9:30pm every night, depending on the sunset.
Singapore
The Singapore Fireworks Celebrations (previously the Singapore Fireworks Festival) is an annual event held in Singapore as part of its National Day celebrations. The festival features local and foreign teams which launch displays on different nights. While currently non-competitive in nature, the organizer has plans to introduce a competitive element in the future.
The annual festival has grown in magnitude, from 4,000 rounds used in 2004, to 6,000 in 2005, to more than 9,100 in 2006.
South Korea
Busan International Fireworks Festival is one of the most significant fireworks festivals in Asia.
Switzerland
In Switzerland fireworks are often used on 1 August, which is a national celebration day.
United Kingdom
One of the biggest occasions for fireworks in the UK is Guy Fawkes Night held each year on 5 November, to celebrate the foiling of the Catholic Gunpowder Plot on 5 November 1605, an attempt to kill King James I. The Guardian newspaper said in 2008 that Britain's biggest Guy Fawkes night events were:
After Dark fireworks, Sheffield
Bangers on the Beach (Holyhead Round Table charity fireworks), Holyhead
Bonfire in Battle, East Sussex
Blackheath Fireworks, London
Bught Park fireworks, Inverness
Fireworks with Vikings, Tutbury, Staffordshire
Flaming Tar Barrels, Ottery St Mary
Glasgow Green fireworks
Halloween Happening fireworks, Derry
Midsummer Common, Cambridge
Sparks in the Park (Cardiff Round Table charity fireworks), Cardiff
The main firework celebrations in the UK are by the public who buy from many suppliers.
United States
America's earliest settlers brought their enthusiasm for fireworks to the United States. Fireworks and black ash were used to celebrate important events long before the American Revolutionary War. The first celebration of Independence Day was in 1777, six years before Americans knew whether or not the new nation would survive the war; fireworks were a part of all festivities. In 1789, George Washington's inauguration was accompanied by a fireworks display.. George Marshall was an American naval hero during the War of 1812 and other campaigns. He was a Master Gunner and pyrotechnics specialist who wrote Marshall's Practical Marine Gunnery in 1822. The book outlines chemical formulas for the composition of fireworks. This early fascination with fireworks' noise and color continues today with fireworks displays commonly included in Independence Day celebrations.
In 2004, Disneyland, in Anaheim, California, pioneered the commercial use of aerial fireworks launched with compressed air rather than gunpowder. The display shell explodes in the air using an electronic timer. The advantages of compressed air launch are a reduction in fumes, and much greater accuracy in height and timing. The Walt Disney Company is now the largest consumer of fireworks in the world.
Halloween
Canada
Fireworks are a popular tradition during Halloween in Vancouver.
Ireland
In the Republic of Ireland and Northern Ireland there are many fireworks displays, during Halloween. The sale of fireworks is strongly restricted in the Republic of Ireland, although many illegal fireworks are sold throughout October or smuggled from Northern Ireland. The maximum punishment for possessing fireworks without a licence, or lighting fireworks in a public place, is a €10,000 fine and a five-year prison sentence.
United States
Two firework displays on All Hallows' Eve in the United States are the "Happy Hallowishes" show at Walt Disney World's Magic Kingdom "Mickey's Not-So-Scary Halloween Party" event, which began in 2005, and the "Halloween Screams" at Disneyland Park, which began in 2009.
Uses other than public displays
In addition to large public displays, people often buy small quantities of fireworks for their own celebrations. Fireworks on general sale are usually less powerful than professional fireworks. Types include firecrackers, rockets, cakes (multishot aerial fireworks), and smoke balls.
Fireworks can also be used in an agricultural capacity as to frighten away birds.
Culture
Competitions
Pyrotechnical competitions are held in many countries. Among them are the Montreal Fireworks Festival, an annual competition held in Montreal, Quebec, Canada; Le Festival d'Art Pyrotechnique, held in the summer annually at the Bay of Cannes in Côte d'Azur, France; and the Philippine International Pyromusical Competition, held in Manila, Philippines amongst the top fireworks companies in the world.
Clubs and organizations
Enthusiasts in the United States have formed clubs which unite hobbyists and professionals. The groups provide safety instruction and organize meetings and private "shoots" at remote premises where members shoot commercial fireworks as well as fire pieces of their own manufacture. Clubs secure permission to fire items otherwise banned by state or local ordinances. Competition among members and between clubs, demonstrating everything from single shells to elaborate displays choreographed to music, are held. One of the oldest clubs is Crackerjacks, Inc., organized in 1976 in the Eastern Seaboard region.
Though based in the US, membership of the Pyrotechnics Guild International, Inc. (PGI) is annual convention founded in 1969, it hosts some the world's biggest fireworks displays occur. Aside from the nightly firework shows, one of the most popular events of the convention is a unique event where individual classes of hand-built fireworks are competitively judged, ranging from simple fireworks rockets to extremely large and complex aerial shells. Some of the biggest, most intricate fireworks displays in the United States take place during the convention week.
Growth
According to industry data, the purchase of fireworks by American consumers has markedly increased since the onset of the COVID-19 pandemic. In the year preceding the pandemic, the consumer fireworks industry reported sales of approximately $1 billion. However, this figure almost doubled in 2020, reaching $2.3 billion. In contrast, commercial fireworks sales amounted to a smaller figure of $400 million. For context, consumer fireworks revenue was significantly lower at $645 million in 2012, indicating a steady growth trend in the market. Industry projections currently anticipate a further increase, forecasting fireworks sales to reach $3.3 billion by 2028.
| Technology | Material and chemical | null |
59497 | https://en.wikipedia.org/wiki/Solubility | Solubility | In chemistry, solubility is the ability of a substance, the solute, to form a solution with another substance, the solvent. Insolubility is the opposite property, the inability of the solute to form such a solution.
The extent of the solubility of a substance in a specific solvent is generally measured as the concentration of the solute in a saturated solution, one in which no more solute can be dissolved. At this point, the two substances are said to be at the solubility equilibrium. For some solutes and solvents, there may be no such limit, in which case the two substances are said to be "miscible in all proportions" (or just "miscible").
The solute can be a solid, a liquid, or a gas, while the solvent is usually solid or liquid. Both may be pure substances, or may themselves be solutions. Gases are always miscible in all proportions, except in very extreme situations, and a solid or liquid can be "dissolved" in a gas only by passing into the gaseous state first.
The solubility mainly depends on the composition of solute and solvent (including their pH and the presence of other dissolved substances) as well as on temperature and pressure. The dependency can often be explained in terms of interactions between the particles (atoms, molecules, or ions) of the two substances, and of thermodynamic concepts such as enthalpy and entropy.
Under certain conditions, the concentration of the solute can exceed its usual solubility limit. The result is a supersaturated solution, which is metastable and will rapidly exclude the excess solute if a suitable nucleation site appears.
The concept of solubility does not apply when there is an irreversible chemical reaction between the two substances, such as the reaction of calcium hydroxide with hydrochloric acid; even though one might say, informally, that one "dissolved" the other. The solubility is also not the same as the rate of solution, which is how fast a solid solute dissolves in a liquid solvent. This property depends on many other variables, such as the physical form of the two substances and the manner and intensity of mixing.
The concept and measure of solubility are extremely important in many sciences besides chemistry, such as geology, biology, physics, and oceanography, as well as in engineering, medicine, agriculture, and even in non-technical activities like painting, cleaning, cooking, and brewing. Most chemical reactions of scientific, industrial, or practical interest only happen after the reagents have been dissolved in a suitable solvent. Water is by far the most common such solvent.
The term "soluble" is sometimes used for materials that can form colloidal suspensions of very fine solid particles in a liquid. The quantitative solubility of such substances is generally not well-defined, however.
Quantification of solubility
The solubility of a specific solute in a specific solvent is generally expressed as the concentration of a saturated solution of the two. Any of the several ways of expressing concentration of solutions can be used, such as the mass, volume, or amount in moles of the solute for a specific mass, volume, or mole amount of the solvent or of the solution.
Per quantity of solvent
In particular, chemical handbooks often express the solubility as grams of solute per 100 millilitres of solvent (g/(100 mL), often written as g/100 ml), or as grams of solute per decilitre of solvent (g/dL); or, less commonly, as grams of solute per litre of solvent (g/L). The quantity of solvent can instead be expressed in mass, as grams of solute per 100 grams of solvent (g/(100 g), often written as g/100 g), or as grams of solute per kilogram of solvent (g/kg). The number may be expressed as a percentage in this case, and the abbreviation "w/w" may be used to indicate "weight per weight". (The values in g/L and g/kg are similar for water, but that may not be the case for other solvents.)
Alternatively, the solubility of a solute can be expressed in moles instead of mass. For example, if the quantity of solvent is given in kilograms, the value is the molality of the solution (mol/kg).
Per quantity of solution
The solubility of a substance in a liquid may also be expressed as the quantity of solute per quantity of solution, rather than of solvent. For example, following the common practice in titration, it may be expressed as moles of solute per litre of solution (mol/L), the molarity of the latter.
In more specialized contexts the solubility may be given by the mole fraction (moles of solute per total moles of solute plus solvent) or by the mass fraction at equilibrium (mass of solute per mass of solute plus solvent). Both are dimensionless numbers between 0 and 1 which may be expressed as percentages (%).
Liquid and gaseous solutes
For solutions of liquids or gases in liquids, the quantities of both substances may be given volume rather than mass or mole amount; such as litre of solute per litre of solvent, or litre of solute per litre of solution. The value may be given as a percentage, and the abbreviation "v/v" for "volume per volume" may be used to indicate this choice.
Conversion of solubility values
Conversion between these various ways of measuring solubility may not be trivial, since it may require knowing the density of the solution — which is often not measured, and cannot be predicted. While the total mass is conserved by dissolution, the final volume may be different from both the volume of the solvent and the sum of the two volumes.
Moreover, many solids (such as acids and salts) will dissociate in non-trivial ways when dissolved; conversely, the solvent may form coordination complexes with the molecules or ions of the solute. In those cases, the sum of the moles of molecules of solute and solvent is not really the total moles of independent particles solution. To sidestep that problem, the solubility per mole of solution is usually computed and quoted as if the solute does not dissociate or form complexes—that is, by pretending that the mole amount of solution is the sum of the mole amounts of the two substances.
Qualifiers used to describe extent of solubility
The extent of solubility ranges widely, from infinitely soluble (without limit, i.e. miscible) such as ethanol in water, to essentially insoluble, such as titanium dioxide in water. A number of other descriptive terms are also used to qualify the extent of solubility for a given application. For example, U.S. Pharmacopoeia gives the following terms, according to the mass msv of solvent required to dissolve one unit of mass msu of solute: (The solubilities of the examples are approximate, for water at 20–25 °C.)
The thresholds to describe something as insoluble, or similar terms, may depend on the application. For example, one source states that substances are described as "insoluble" when their solubility is less than 0.1 g per 100 mL of solvent.
Molecular view
Solubility occurs under dynamic equilibrium, which means that solubility results from the simultaneous and opposing processes of dissolution and phase joining (e.g. precipitation of solids). A stable state of the solubility equilibrium occurs when the rates of dissolution and re-joining are equal, meaning the relative amounts of dissolved and non-dissolved materials are equal. If the solvent is removed, all of the substance that had dissolved is recovered.
The term solubility is also used in some fields where the solute is altered by solvolysis. For example, many metals and their oxides are said to be "soluble in hydrochloric acid", although in fact the aqueous acid irreversibly degrades the solid to give soluble products. Most ionic solids dissociate when dissolved in polar solvents. In those cases where the solute is not recovered upon evaporation of the solvent, the process is referred to as solvolysis. The thermodynamic concept of solubility does not apply straightforwardly to solvolysis.
When a solute dissolves, it may form several species in the solution. For example, an aqueous solution of cobalt(II) chloride can afford , each of which interconverts.
Factors affecting solubility
Solubility is defined for specific phases. For example, the solubility of aragonite and calcite in water are expected to differ, even though they are both polymorphs of calcium carbonate and have the same chemical formula.
The solubility of one substance in another is determined by the balance of intermolecular forces between the solvent and solute, and the entropy change that accompanies the solvation. Factors such as temperature and pressure will alter this balance, thus changing the solubility.
Solubility may also strongly depend on the presence of other species dissolved in the solvent, for example, complex-forming anions (ligands) in liquids. Solubility will also depend on the excess or deficiency of a common ion in the solution, a phenomenon known as the common-ion effect. To a lesser extent, solubility will depend on the ionic strength of solutions. The last two effects can be quantified using the equation for solubility equilibrium.
For a solid that dissolves in a redox reaction, solubility is expected to depend on the potential (within the range of potentials under which the solid remains the thermodynamically stable phase). For example, solubility of gold in high-temperature water is observed to be almost an order of magnitude higher (i.e. about ten times higher) when the redox potential is controlled using a highly oxidizing Fe3O4-Fe2O3 redox buffer than with a moderately oxidizing Ni-NiO buffer.
Solubility (metastable, at concentrations approaching saturation) also depends on the physical size of the crystal or droplet of solute (or, strictly speaking, on the specific surface area or molar surface area of the solute). For quantification, see the equation in the article on solubility equilibrium. For highly defective crystals, solubility may increase with the increasing degree of disorder. Both of these effects occur because of the dependence of solubility constant on the Gibbs energy of the crystal. The last two effects, although often difficult to measure, are of practical importance. For example, they provide the driving force for precipitate aging (the crystal size spontaneously increasing with time).
Temperature
The solubility of a given solute in a given solvent is function of temperature. Depending on the change in enthalpy (ΔH) of the dissolution reaction, i.e., on the endothermic (ΔH > 0) or exothermic (ΔH < 0) character of the dissolution reaction, the solubility of a given compound may increase or decrease with temperature. The van 't Hoff equation relates the change of solubility equilibrium constant (Ksp) to temperature change and to reaction enthalpy change. For most solids and liquids, their solubility increases with temperature because their dissolution reaction is endothermic (ΔH > 0). In liquid water at high temperatures, (e.g. that approaching the critical temperature), the solubility of ionic solutes tends to decrease due to the change of properties and structure of liquid water; the lower dielectric constant results in a less polar solvent and in a change of hydration energy affecting the ΔG of the dissolution reaction.
Gaseous solutes exhibit more complex behavior with temperature. As the temperature is raised, gases usually become less soluble in water (exothermic dissolution reaction related to their hydration) (to a minimum, which is below 120 °C for most permanent gases), but more soluble in organic solvents (endothermic dissolution reaction related to their solvation).
The chart shows solubility curves for some typical solid inorganic salts in liquid water (temperature is in degrees Celsius, i.e. kelvins minus 273.15). Many salts behave like barium nitrate and disodium hydrogen arsenate, and show a large increase in solubility with temperature (ΔH > 0). Some solutes (e.g. sodium chloride in water) exhibit solubility that is fairly independent of temperature (ΔH ≈ 0). A few, such as calcium sulfate (gypsum) and cerium(III) sulfate, become less soluble in water as temperature increases (ΔH < 0). This is also the case for calcium hydroxide (portlandite), whose solubility at 70 °C is about half of its value at 25 °C. The dissolution of calcium hydroxide in water is also an exothermic process (ΔH < 0). As dictated by the van 't Hoff equation and Le Chatelier's principle, low temperatures favor dissolution of Ca(OH)2. Portlandite solubility increases at low temperature. This temperature dependence is sometimes referred to as "retrograde" or "inverse" solubility. Occasionally, a more complex pattern is observed, as with sodium sulfate, where the less soluble decahydrate crystal (mirabilite) loses water of crystallization at 32 °C to form a more soluble anhydrous phase (thenardite) with a smaller change in Gibbs free energy (ΔG) in the dissolution reaction.
The solubility of organic compounds nearly always increases with temperature. The technique of recrystallization, used for purification of solids, depends on a solute's different solubilities in hot and cold solvent. A few exceptions exist, such as certain cyclodextrins.
Pressure
For condensed phases (solids and liquids), the pressure dependence of solubility is typically weak and usually neglected in practice. Assuming an ideal solution, the dependence can be quantified as:
where the index iterates the components, is the mole fraction of the -th component in the solution, is the pressure, the index refers to constant temperature, is the partial molar volume of the -th component in the solution, is the partial molar volume of the -th component in the dissolving solid, and is the universal gas constant.
The pressure dependence of solubility does occasionally have practical significance. For example, precipitation fouling of oil fields and wells by calcium sulfate (which decreases its solubility with decreasing pressure) can result in decreased productivity with time.
Solubility of gases
Henry's law is used to quantify the solubility of gases in solvents. The solubility of a gas in a solvent is directly proportional to the partial pressure of that gas above the solvent. This relationship is similar to Raoult's law and can be written as:
where is a temperature-dependent constant (for example, 769.2 L·atm/mol for dioxygen (O2) in water at 298 K), is the partial pressure (in atm), and is the concentration of the dissolved gas in the liquid (in mol/L).
The solubility of gases is sometimes also quantified using Bunsen solubility coefficient.
In the presence of small bubbles, the solubility of the gas does not depend on the bubble radius in any other way than through the effect of the radius on pressure (i.e. the solubility of gas in the liquid in contact with small bubbles is increased due to pressure increase by Δp = 2γ/r; see Young–Laplace equation).
Henry's law is valid for gases that do not undergo change of chemical speciation on dissolution. Sieverts' law shows a case when this assumption does not hold.
The carbon dioxide solubility in seawater is also affected by temperature, pH of the solution, and by the carbonate buffer. The decrease of solubility of carbon dioxide in seawater when temperature increases is also an important retroaction factor (positive feedback) exacerbating past and future climate changes as observed in ice cores from the Vostok site in Antarctica. At the geological time scale, because of the Milankovich cycles, when the astronomical parameters of the Earth orbit and its rotation axis progressively change and modify the solar irradiance at the Earth surface, temperature starts to increase. When a deglaciation period is initiated, the progressive warming of the oceans releases CO2 into the atmosphere because of its lower solubility in warmer sea water. In turn, higher levels of CO2 in the atmosphere increase the greenhouse effect and carbon dioxide acts as an amplifier of the general warming.
Polarity
A popular aphorism used for predicting solubility is "like dissolves like" also expressed in the Latin language as "Similia similibus solventur". This statement indicates that a solute will dissolve best in a solvent that has a similar chemical structure to itself, based on favorable entropy of mixing. This view is simplistic, but it is a useful rule of thumb. The overall solvation capacity of a solvent depends primarily on its polarity. For example, a very polar (hydrophilic) solute such as urea is very soluble in highly polar water, less soluble in fairly polar methanol, and practically insoluble in non-polar solvents such as benzene. In contrast, a non-polar or lipophilic solute such as naphthalene is insoluble in water, fairly soluble in methanol, and highly soluble in non-polar benzene.
In even more simple terms a simple ionic compound (with positive and negative ions) such as sodium chloride (common salt) is easily soluble in a highly polar solvent (with some separation of positive (δ+) and negative (δ-) charges in the covalent molecule) such as water, as thus the sea is salty as it accumulates dissolved salts since early geological ages.
The solubility is favored by entropy of mixing (ΔS) and depends on enthalpy of dissolution (ΔH) and the hydrophobic effect. The free energy of dissolution (Gibbs energy) depends on temperature and is given by the relationship: ΔG = ΔH – TΔS. Smaller ΔG means greater solubility.
Chemists often exploit differences in solubilities to separate and purify compounds from reaction mixtures, using the technique of liquid-liquid extraction. This applies in vast areas of chemistry from drug synthesis to spent nuclear fuel reprocessing.
Rate of dissolution
Dissolution is not an instantaneous process. The rate of solubilization (in kg/s) is related to the solubility product and the surface area of the material. The speed at which a solid dissolves may depend on its crystallinity or lack thereof in the case of amorphous solids and the surface area (crystallite size) and the presence of polymorphism. Many practical systems illustrate this effect, for example in designing methods for controlled drug delivery. In some cases, solubility equilibria can take a long time to establish (hours, days, months, or many years; depending on the nature of the solute and other factors).
The rate of dissolution can be often expressed by the Noyes–Whitney equation or the Nernst and Brunner equation of the form:
where:
= mass of dissolved material
= time
= surface area of the interface between the dissolving substance and the solvent
= diffusion coefficient
= thickness of the boundary layer of the solvent at the surface of the dissolving substance
= mass concentration of the substance on the surface
= mass concentration of the substance in the bulk of the solvent
For dissolution limited by diffusion (or mass transfer if mixing is present), is equal to the solubility of the substance. When the dissolution rate of a pure substance is normalized to the surface area of the solid (which usually changes with time during the dissolution process), then it is expressed in kg/m2s and referred to as "intrinsic dissolution rate". The intrinsic dissolution rate is defined by the United States Pharmacopeia.
Dissolution rates vary by orders of magnitude between different systems. Typically, very low dissolution rates parallel low solubilities, and substances with high solubilities exhibit high dissolution rates, as suggested by the Noyes-Whitney equation.
Theories of solubility
Solubility product
Solubility constants are used to describe saturated solutions of ionic compounds of relatively low solubility (see solubility equilibrium). The solubility constant is a special case of an equilibrium constant. Since it is a product of ion concentrations in equilibrium, it is also known as the solubility product. It describes the balance between dissolved ions from the salt and undissolved salt. The solubility constant is also "applicable" (i.e. useful) to precipitation, the reverse of the dissolving reaction. As with other equilibrium constants, temperature can affect the numerical value of solubility constant. While the solubility constant is not as simple as solubility, the value of this constant is generally independent of the presence of other species in the solvent.
Other theories
The Flory–Huggins solution theory is a theoretical model describing the solubility of polymers. The Hansen solubility parameters and the Hildebrand solubility parameters are empirical methods for the prediction of solubility. It is also possible to predict solubility from other physical constants such as the enthalpy of fusion.
The octanol-water partition coefficient, usually expressed as its logarithm (Log P), is a measure of differential solubility of a compound in a hydrophobic solvent (1-octanol) and a hydrophilic solvent (water). The logarithm of these two values enables compounds to be ranked in terms of hydrophilicity (or hydrophobicity).
The energy change associated with dissolving is usually given per mole of solute as the enthalpy of solution.
Applications
Solubility is of fundamental importance in a large number of scientific disciplines and practical applications, ranging from ore processing and nuclear reprocessing to the use of medicines, and the transport of pollutants.
Solubility is often said to be one of the "characteristic properties of a substance", which means that solubility is commonly used to describe the substance, to indicate a substance's polarity, to help to distinguish it from other substances, and as a guide to applications of the substance. For example, indigo is described as "insoluble in water, alcohol, or ether but soluble in chloroform, nitrobenzene, or concentrated sulfuric acid".
Solubility of a substance is useful when separating mixtures. For example, a mixture of salt (sodium chloride) and silica may be separated by dissolving the salt in water, and filtering off the undissolved silica. The synthesis of chemical compounds, by the milligram in a laboratory, or by the ton in industry, both make use of the relative solubilities of the desired product, as well as unreacted starting materials, byproducts, and side products to achieve separation.
Another example of this is the synthesis of benzoic acid from phenylmagnesium bromide and dry ice. Benzoic acid is more soluble in an organic solvent such as dichloromethane or diethyl ether, and when shaken with this organic solvent in a separatory funnel, will preferentially dissolve in the organic layer. The other reaction products, including the magnesium bromide, will remain in the aqueous layer, clearly showing that separation based on solubility is achieved. This process, known as liquid–liquid extraction, is an important technique in synthetic chemistry. Recycling is used to ensure maximum extraction.
Differential solubility
In flowing systems, differences in solubility often determine the dissolution-precipitation driven transport of species. This happens when different parts of the system experience different conditions. Even slightly different conditions can result in significant effects, given sufficient time.
For example, relatively low solubility compounds are found to be soluble in more extreme environments, resulting in geochemical and geological effects of the activity of hydrothermal fluids in the Earth's crust. These are often the source of high quality economic mineral deposits and precious or semi-precious gems. In the same way, compounds with low solubility will dissolve over extended time (geological time), resulting in significant effects such as extensive cave systems or Karstic land surfaces.
Solubility of ionic compounds in water
Some ionic compounds (salts) dissolve in water, which arises because of the attraction between positive and negative charges (see: solvation). For example, the salt's positive ions (e.g. Ag+) attract the partially negative oxygen atom in . Likewise, the salt's negative ions (e.g. Cl−) attract the partially positive hydrogens in . Note: the oxygen atom is partially negative because it is more electronegative than hydrogen, and vice versa (see: chemical polarity).
However, there is a limit to how much salt can be dissolved in a given volume of water. This concentration is the solubility and related to the solubility product, Ksp. This equilibrium constant depends on the type of salt ( vs. , for example), temperature, and the common ion effect.
One can calculate the amount of that will dissolve in 1 liter of pure water as follows:
Ksp = [Ag+] × [Cl−] / M2 (definition of solubility product; M = mol/L)
Ksp = 1.8 × 10−10 (from a table of solubility products)
[Ag+] = [Cl−], in the absence of other silver or chloride salts, so
[Ag+]2 = 1.8 × 10−10 M2
[Ag+] = 1.34 × 10−5 mol/L
The result: 1 liter of water can dissolve 1.34 × 10−5 moles of at room temperature. Compared with other salts, is poorly soluble in water. For instance, table salt () has a much higher Ksp = 36 and is, therefore, more soluble. The following table gives an overview of solubility rules for various ionic compounds.
Solubility of organic compounds
The principle outlined above under polarity, that like dissolves like, is the usual guide to solubility with organic systems. For example, petroleum jelly will dissolve in gasoline because both petroleum jelly and gasoline are non-polar hydrocarbons. It will not, on the other hand, dissolve in ethyl alcohol or water, since the polarity of these solvents is too high. Sugar will not dissolve in gasoline, since sugar is too polar in comparison with gasoline. A mixture of gasoline and sugar can therefore be separated by filtration or extraction with water.
Solid solution
This term is often used in the field of metallurgy to refer to the extent that an alloying element will dissolve into the base metal without forming a separate phase. The solvus or solubility line (or curve) is the line (or lines) on a phase diagram that give the limits of solute addition. That is, the lines show the maximum amount of a component that can be added to another component and still be in solid solution. In the solid's crystalline structure, the 'solute' element can either take the place of the matrix within the lattice (a substitutional position; for example, chromium in iron) or take a place in a space between the lattice points (an interstitial position; for example, carbon in iron).
In microelectronic fabrication, solid solubility refers to the maximum concentration of impurities one can place into the substrate.
In solid compounds (as opposed to elements), the solubility of a solute element can also depend on the phases separating out in equilibrium. For example, amount of Sn soluble in the ZnSb phase can depend significantly on whether the phases separating out in equilibrium are (Zn4Sb3+Sn(L)) or (ZnSnSb2+Sn(L)). Besides these, the ZnSb compound with Sn as a solute can separate out into other combinations of phases after the solubility limit is reached depending on the initial chemical composition during synthesis. Each combination produces a different solubility of Sn in ZnSb. Hence solubility studies in compounds, concluded upon the first instance of observing secondary phases separating out might underestimate solubility. While the maximum number of phases separating out at once in equilibrium can be determined by the Gibb's phase rule, for chemical compounds there is no limit on the number of such phase separating combinations itself. Hence, establishing the "maximum solubility" in solid compounds experimentally can be difficult, requiring equilibration of many samples. If the dominant crystallographic defect (mostly interstitial or substitutional point defects) involved in the solid-solution can be chemically intuited beforehand, then using some simple thermodynamic guidelines can considerably reduce the number of samples required to establish maximum solubility.
Incongruent dissolution
Many substances dissolve congruently (i.e. the composition of the solid and the dissolved solute stoichiometrically match). However, some substances may dissolve incongruently, whereby the composition of the solute in solution does not match that of the solid. This solubilization is accompanied by alteration of the "primary solid" and possibly formation of a secondary solid phase. However, in general, some primary solid also remains and a complex solubility equilibrium establishes. For example, dissolution of albite may result in formation of gibbsite.
.
In this case, the solubility of albite is expected to depend on the solid-to-solvent ratio. This kind of solubility is of great importance in geology, where it results in formation of metamorphic rocks.
In principle, both congruent and incongruent dissolution can lead to the formation of secondary solid phases in equilibrium. So, in the field of Materials Science, the solubility for both cases is described more generally on chemical composition phase diagrams.
Solubility prediction
Solubility is a property of interest in many aspects of science, including but not limited to: environmental predictions, biochemistry, pharmacy, drug-design, agrochemical design, and protein ligand binding. Aqueous solubility is of fundamental interest owing to the vital biological and transportation functions played by water. In addition, to this clear scientific interest in water solubility and solvent effects; accurate predictions of solubility are important industrially. The ability to accurately predict a molecule's solubility represents potentially large financial savings in many chemical product development processes, such as pharmaceuticals. In the pharmaceutical industry, solubility predictions form part of the early stage lead optimisation process of drug candidates. Solubility remains a concern all the way to formulation. A number of methods have been applied to such predictions including quantitative structure–activity relationships (QSAR), quantitative structure–property relationships (QSPR) and data mining. These models provide efficient predictions of solubility and represent the current standard. The draw back such models is that they can lack physical insight. A method founded in physical theory, capable of achieving similar levels of accuracy at an sensible cost, would be a powerful tool scientifically and industrially.
Methods founded in physical theory tend to use thermodynamic cycles, a concept from classical thermodynamics. The two common thermodynamic cycles used involve either the calculation of the free energy of sublimation (solid to gas without going through a liquid state) and the free energy of solvating a gaseous molecule (gas to solution), or the free energy of fusion (solid to a molten phase) and the free energy of mixing (molten to solution). These two process are represented in the following diagrams.
These cycles have been used for attempts at first principles predictions (solving using the fundamental physical equations) using physically motivated solvent models, to create parametric equations and QSPR models and combinations of the two. The use of these cycles enables the calculation of the solvation free energy indirectly via either gas (in the sublimation cycle) or a melt (fusion cycle). This is helpful as calculating the free energy of solvation directly is extremely difficult. The free energy of solvation can be converted to a solubility value using various formulae, the most general case being shown below, where the numerator is the free energy of solvation, R is the gas constant and T is the temperature in kelvins.
Well known fitted equations for solubility prediction are the general solubility equations. These equations stem from the work of Yalkowsky et al. The original formula is given first, followed by a revised formula which takes a different assumption of complete miscibility in octanol.
These equations are founded on the principles of the fusion cycle.
| Physical sciences | Mixture | Chemistry |
59597 | https://en.wikipedia.org/wiki/Disc%20brake | Disc brake | A disc brake is a type of brake that uses the calipers to squeeze pairs of pads against a disc (sometimes called a [brake] rotor) to create friction. There are two basic types of brake pad friction mechanisms: abrasive friction and adherent friction. This action slows the rotation of a shaft, such as a vehicle axle, either to reduce its rotational speed or to hold it stationary. The energy of motion is converted into heat, which must be dissipated to the environment.
Hydraulically actuated disc brakes are the most commonly used mechanical device for slowing motor vehicles. The principles of a disc brake apply to almost any rotating shaft. The components include the disc, master cylinder, and caliper, which contain at least one cylinder and two brake pads on both sides of the rotating disc.
Design
The development of disc-type brakes began in England in the 1890s. In 1902, the Lanchester Motor Company designed brakes that looked and operated similarly to a modern disc-brake system even though the disc was thin and a cable activated the brake pad. Other designs were not practical or widely available in cars for another 60 years. Successful application began in airplanes before World War II. The German Tiger tank was fitted with discs in 1942. After the war, technological progress began in 1949, with caliper-type four-wheel disc brakes on the Crosley line and a Chrysler non-caliper type. In the 1950s, there was a demonstration of superiority at the 1953 24 Hours of Le Mans race, which required braking from high speeds several times per lap. The Jaguar racing team won, using disc brake-equipped cars, with much of the credit being given to the brakes' superior performance over rivals equipped with drum brakes. Mass production began with the 1949–1950 inclusion in all Crosley production, with sustained mass production starting in 1955 Citroën DS.
Disc brakes offer better stopping performance than drum brakes because the disc is more readily cooled. Consequently, discs are less prone to the brake fade caused when brake components overheat. Disc brakes also recover more quickly from immersion (wet brakes are less effective than dry ones).
Most drum brake designs have at least one leading shoe, which gives a servo-effect. By contrast, a disc brake has no self-servo effect, and its braking force is always proportional to the pressure placed on the brake pad by the braking system via any brake servo, brake pedal, or lever. This tends to give the driver a better "feel" and helps to avoid impending lockup. Drums are also prone to "bell mouthing" and trap worn lining material within the assembly, causing various braking problems.
The disc is usually made of cast iron. In some cases, it may be made of composites such as reinforced carbon–carbon or ceramic matrix composites. This is connected to the wheel and the axle. To slow down the wheel, friction material in the form of brake pads, mounted on the brake caliper, is forced mechanically, hydraulically, pneumatically, or electromagnetically against both sides of the disc. Friction causes the disc and attached wheel to slow or stop.
Operation
The brake disc is the rotating part of a wheel's disc brake assembly, against which the brake pads are applied. The material is typically gray iron, a form of cast iron. The design of the discs varies. Some are solid, but others are hollowed out with fins or vanes joining the disc's two contact surfaces (usually included in the casting process). The weight and power of the vehicle determine the need for ventilated discs. The "ventilated" disc design helps to dissipate the generated heat and is commonly used on the more heavily loaded front discs.
Discs for motorcycles, bicycles, and many cars often have holes or slots cut through the disc. This is done for better heat dissipation, to aid surface-water dispersal, to reduce noise, to reduce mass, or purely for non-functional aesthetics.
Slotted discs have shallow channels machined into the disc to aid in removing dust and gas. Slotting is preferred in most racing environments to remove gas and water and deglaze brake pads. Some discs are both drilled and slotted. Slotted discs are generally not used on standard vehicles because they quickly wear down brake pads; however, removing of material is beneficial to race vehicles since it keeps the pads soft and avoids vitrification of their surfaces. On the road, drilled or slotted discs still have a positive effect in wet conditions because the holes or slots prevent a film of water from building up between the disc and the pads.
Two-piece discs consist of a central section combined with a separately manufactured outer friction ring. The central section is often called a bell or hat because of its shape. It is commonly manufactured from an alloy such as a 7075 alloy and hard anodised for a lasting finish. The outer disc ring is usually made of grey iron. They can also be made of steel or carbon ceramic for particular applications. These materials originated from motorsport use and are available in high-performance vehicles and aftermarket upgrades. Two-piece discs can be supplied as a fixed assembly with regular nuts, bolts, and washers or a more complicated floating system where drive bobbins allow the two parts of the brake disc to expand and contract at different rates, therefore reducing the chance of a disc warping from overheating. Key advantages of a two-piece disc are a reduction of critical un-sprung weight and the dissipation of heat from the disc surface through the alloy bell (hat). Both fixed and floating options have their drawbacks and advantages. Floating discs are prone to rattle and collection of debris and are best suited to motorsport, whereas fixed discs are best for road use.
History
Early experiments
The development of disc brakes began in England in the 1890s. The first caliper-type automobile disc brake was patented by Frederick William Lanchester in his Birmingham factory in 1902 and used successfully on Lanchester cars. However, the limited choice of metals in this period meant he used copper as the braking medium acting on the disc. The poor state of the roads at this time, no more than dusty, rough tracks, meant that the copper wore quickly, making the system impractical.
In 1921, the Douglas motorcycle company introduced a form of disc brake on the front wheel of their overhead-valve sports models. Patented by the British Motorcycle & Cycle-Car Research Association, Douglas described the device as a "novel wedge brake" working on a "beveled hub flange". A Bowden cable operated the brake. Front and rear brakes of this type were fitted to the machine on which Tom Sheard rode to victory in the 1923 Senior TT.
Successful application began on railroad streamliner passenger trains, airplanes, and tanks before and during World War II. In the US, the Budd Company introduced disc brakes on the General Pershing Zephyr for the Burlington Railroad in 1938. By the early 1950s, disc brakes were regularly applied to new passenger rolling stock. In Britain, the Daimler Company used disc brakes on its Daimler Armoured Car of 1939. The disc brakes, made by the Girling company, were necessary because in that four-wheel drive (4×4) vehicle the epicyclic final drive was in the wheel hubs and therefore left no room for conventional hub-mounted drum brakes.
At Germany's Argus Motoren, Hermann Klaue (1912-2001) had patented disc brakes in 1940. Argus supplied wheels fitted with disc brakes e.g. for the Arado Ar 96. The German Tiger I heavy tank, was introduced in 1942 with a 55 cm Argus-Werke disc on each drive shaft.
The American Crosley Hot Shot had four-wheel disc brakes in 1949 and 1950. However, these quickly proved troublesome and were removed. Crosley returned to drum brakes, and drum brake conversions for Hot Shots were popular. Lack of sufficient research caused reliability problems, such as sticking and corrosion, especially in regions using salt on winter roads. Crosley four-wheel disc brakes made the cars, and Crosley-based specials, popular in SCCA H-Production and H-modified racing in the 1950s. The Crosley disc was a Goodyear-Hawley design, a modern caliper "spot" type with a modern disc, derived from a design from aircraft applications.
Chrysler developed a unique braking system, offered from 1949 until 1953. Instead of the disc with caliper squeezing on it, this system used twin expanding discs that rubbed against the inner surface of a cast-iron brake drum, which doubled as the brake housing. The discs spread apart to create friction against the inner drum surface through the action of standard wheel cylinders. Because of the expense, the brakes were only standard on the Chrysler Crown and the Town and Country Newport in 1950. They were optional, however, on other Chryslers, priced around $400, at a time when an entire Crosley Hot Shot retailed for $935. This four-wheel disc brake system was built by Auto Specialties Manufacturing Company (Ausco) of St. Joseph, Michigan, under patents of inventor H.L. Lambert, and was first tested on a 1939 Plymouth. Chrysler discs were "self-energizing," in that some of the braking energy itself contributed to the braking effort. This was accomplished by small balls set into oval holes leading to the braking surface. When the disc made initial contact with the friction surface, the balls would be forced up the holes forcing the discs further apart and augmenting the braking energy. This made for lighter braking pressure than with calipers, avoided brake fade, promoted cooler running, and provided one-third more friction surface than standard Chrysler twelve-inch drums. Today's owners consider the Ausco-Lambert very reliable and powerful, but admit its grabbiness and sensitivity.
In 1953, 50 aluminum-bodied Austin-Healey 100S (Sebring) models, built primarily for racing, were the first European cars sold to the public to have disc brakes, fitted to all four wheels.
First impact in racing
The Jaguar C-Type racing car won the 1953 24 Hours of Le Mans, the only vehicle in the race to use disc brakes, developed in the UK by Dunlop, and the first car at Le Mans ever to average over 100 mph. "Rivals' large drum brakes could match discs' ultimate stopping, but not their formidable staying power."
Before this, in 1950, a Crosley HotShot with stock four-wheel disc brakes won the Index of Performance in the first race at Sebring (six hours rather than 12) on New Year's Eve in 1950.
Mass production
The Citroën DS was the first sustained mass production use of modern automotive disc brakes, in 1955. The car featured caliper-type front disc brakes among its many innovations. These discs were mounted inboard near the transmission and were powered by the vehicle's central hydraulic system. This model went on to sell 1.5 million units over 20 years with the same brake setup.
Despite early experiments in 1902, from British Lanchester Motor Company, and in 1949 from Americans Chrysler and Crosley, the costly, trouble-prone technology was not ready for mass production. Attempts were soon withdrawn.
The Jensen 541, with four-wheel disc brakes, followed in 1956. Triumph exhibited a 1956 TR3 with disc brakes to the public, but the first production cars with Girling front-disc brakes were made in September 1956. Jaguar began to offer disc brakes in February 1957 on the XK150 model, soon to follow with the Mark 1 sports saloon and in 1959 with the Mark IX large saloon.
Disc brakes were most popular on sports cars when they were first introduced since these vehicles are more demanding about brake performance. Discs have now become the more common form in most passenger vehicles. However, many (lightweight vehicles) use drum brakes on the rear wheels to keep costs and weight down as well as to simplify the provisions for a parking brake. This can be a reasonable compromise because the front brakes perform most of the braking effort.
Many early implementations for automobiles located the brakes on the inboard side of the driveshaft, near the differential, while most brakes today are located inside the wheels. An inboard location reduces the unsprung weight and eliminates a source of heat transfer to the tires.
Historically, brake discs were manufactured worldwide with a concentration in Europe and America. Between 1989 and 2005, the manufacturing of brake discs migrated predominantly to China.
In the U.S.
In 1963, the Studebaker Avanti was factory-equipped with front disc brakes as standard equipment. This Bendix system licensed from Dunlop was also optional on some of the other Studebaker models. Front disc brakes became standard equipment on the 1965 Rambler Marlin. The Bendix units were optional on all American Motors' Rambler Classic and Ambassador models as well as on the Ford Thunderbird, and the Lincoln Continental. A four-wheel disc brake system was also introduced in 1965 on the Chevrolet Corvette Stingray. Most U.S. cars switched from front drum brakes to front disk brakes in the late 1970s and early 1980s.
Motorcycles and scooters
Lambretta introduced the first high-volume production use of a single, floating, front disc brake, enclosed in a ventilated cast alloy hub and actuated by cable, on the 1962 TV175. This was followed by the GT200 in 1964.
MV Agusta was the second manufacturer to offer a front disc brake motorcycle to the public on a small scale in 1965, on their expensive 600 touring motorcycle featuring cable-operated mechanical actuation. In 1969, Honda introduced the more affordable CB750, which had a single hydraulically actuated front disc brake (and a rear drum brake), and which sold in huge numbers.
Unlike cars, disc brakes that are located within the wheel, bike disc brakes are in the airstream and have optimum cooling. Although cast iron discs have a porous surface that provides superior braking performance, such discs rust in the rain and become unsightly. Accordingly, motorcycle discs are usually stainless steel, drilled, slotted, or wavy to disperse rainwater. Modern motorcycle discs tend to have a floating design whereby the disc "floats" on bobbins and can move slightly, allowing better disc centering with a fixed caliper. A floating disc also avoids disc warping and reduces heat transfer to the wheel hub.
Calipers have evolved from simple single-piston units to two-, four- and even six-piston items. Compared to cars, motorcycles have a higher center of mass:wheelbase ratio, so they experience more weight transfer when braking. Front brakes absorb most of the braking forces, while the rear brake serves mainly to balance the motorcycle during braking. Modern sport bikes typically have twin large front discs, with a much smaller single rear disc. Bikes that are particularly fast or heavy may have vented discs.
Early disc brakes (such as on the early Honda Fours and the Norton Commando) sited the calipers on top of the disc, ahead of the fork slider. Although this gave the brake pads better cooling, it is now almost universal practice to site the caliper behind the slider (to reduce the angular momentum of the fork assembly). Rear disc calipers may be mounted above (e.g. BMW R1100S) or below (e.g. Yamaha TRX850) the swinging arm: a low mount provides for a marginally lower center of gravity, while an upper siting keeps the caliper cleaner and better-protected from road obstacles.
One problem with motorcycle disc brakes is that when a bike gets into a violent tank-slapper (high-speed oscillation of the front wheel) the brake pads in the calipers are forced away from the discs, so when the rider applies the brake lever, the caliper pistons push the pads towards the discs without actually making contact. The rider then brakes harder, forcing the pads onto the disc much more aggressively than standard braking. An example of this was the Michele Pirro incident at Mugello, Italy 1 June 2018. At least one manufacturer has developed a system to counter the pads being forced away.
A modern development, particularly on inverted ("upside down", or "USD") forks is the radially mounted caliper. Although these are fashionable, there is no evidence that they improve braking performance or add to the fork's stiffness. (Lacking the option of a fork brace, USD forks may be best stiffened by an oversized front axle).
Bicycles
Bike disc brakes may range from simple, mechanical (cable) systems, to expensive and powerful, multi-piston hydraulic disc systems, commonly used on downhill racing bikes. Improved technology has seen the creation of vented discs for use on mountain bikes, similar to those on cars, introduced to help avoid heat fade on fast alpine descents. Discs are also used on road bicycles for all-weather cycling with predictable braking. By 2024, almost all road bikes are equipped with disc brakes, just like Mountain bikes. Drums are sometimes preferred as harder to damage in crowded parking, where discs are sometimes bent. Most bicycle brake discs are made of steel. Stainless steel is preferred due to its anti-rust properties. Discs are thin, often about 2 mm. Some use a two-piece floating disc style, others use a one-piece solid metal disc. Bicycle disc brakes use either a two-piston caliper that clamps the disc from both sides or a single-piston caliper with one moving pad that contacts the disc first, and then pushes the disc against the non-moving pad. Because energy efficiency is so important in bicycles, an uncommon feature of bicycle brakes is that the pads retract to eliminate residual drag when the brake is released. In contrast, most other brakes drag the pads lightly when released to minimize initial operational travel.
Heavy vehicles
Disc brakes are increasingly used on very large and heavy road vehicles, where previously large drum brakes were nearly universal. One reason is that the disc's lack of self-assist makes brake force much more predictable, so peak brake force can be raised without more risk of braking-induced steering or jackknifing on articulated vehicles. Another is disc brakes fade less when hot, and in a heavy vehicle air and rolling drag and engine braking are small parts of total braking force, so brakes are used harder than on lighter vehicles, and drum brake fade can occur in a single stop. For these reasons, a heavy truck with disc brakes can stop in about 120% of the distance of a passenger car, but with drums, stopping takes about 150% of the distance. In Europe, stopping distance regulations essentially require disc brakes for heavy vehicles. In the U.S., drums are allowed and are typically preferred for their lower purchase price, despite higher total lifetime cost and more frequent service intervals.
Rail and aircraft
Still-larger discs are used for railroad cars, trams, and some airplanes. Passenger rail cars and light rail vehicles often use disc brakes outboard of the wheels, which helps ensure a free flow of cooling air. Some modern passenger rail cars, such as the Amfleet II cars, use inboard disc brakes. This reduces wear from debris and provides protection from rain and snow, which would make the discs slippery and unreliable. However, there is still plenty of cooling for reliable operation. Some airplanes have the brake mounted with very little cooling, and the brake gets hot when stopping. This is acceptable as there is sufficient time for cooling, where the maximum braking energy is very predictable. Should the braking energy exceed the maximum, for example during an emergency occurring during take-off, aircraft wheels can be fitted with a fusible plug to prevent the tire bursting. This is a milestone test in aircraft development.
Automotive use
For automotive use, disc brake discs are commonly made of grey iron. The SAE maintains a specification for the manufacture of grey iron for various applications. For normal car and light-truck applications, SAE specification J431 G3000 (superseded to G10) dictates the correct range of hardness, chemical composition, tensile strength, and other properties necessary for the intended use. Some racing cars and airplanes use brakes with carbon fiber discs and carbon fiber pads to reduce weight. Wear rates tend to be high, and braking may be poor or grabby until the brake is hot.
Racing
In racing and high-performance road cars, other disc materials have been employed. Reinforced carbon discs and pads inspired by aircraft braking systems such as those used on Concorde were introduced in Formula One by Brabham in conjunction with Dunlop in 1976. Carbon–carbon braking is now used in most top-level motorsport worldwide, reducing unsprung weight, giving better frictional performance and improved structural properties at high temperatures, compared to cast iron. Carbon brakes have occasionally been applied to road cars, by the French Venturi sports car manufacturer in the mid-1990s for example, but need to reach a very high operating temperature before becoming truly effective and so are not well suited to road use. The extreme heat generated in these systems is visible during night racing, especially on shorter tracks. It is not uncommon to see the brake discs glowing red during use.
Ceramic composites
Ceramic discs are used in some high-performance cars and heavy vehicles.
The first development of the modern ceramic brake was made by British engineers for TGV applications in 1988. The objective was to reduce weight, and the number of brakes per axle, as well as provide stable friction from high speeds and all temperatures. The result was a carbon-fiber-reinforced ceramic process which is now used in various forms for automotive, railway, and aircraft brake applications.
Due to the high heat tolerance and mechanical strength of ceramic composite discs, they are often used on exotic vehicles where the cost is not prohibitive. They are also found in industrial applications where the ceramic disc's lightweight and low-maintenance properties justify the cost. Composite brakes can withstand temperatures that would damage steel discs.
Porsche's Composite Ceramic Brakes (PCCB) are siliconized carbon fiber, with high-temperature capability, a 50% weight reduction over iron discs (hence reducing the vehicle's unsprung weight), a significant reduction in dust generation, substantially extended maintenance intervals, and enhanced durability in corrosive environments. Found on some of their more expensive models, it is also an optional brake for all street Porsches at added expense. They can be recognized by the bright yellow paintwork on the aluminum six-piston calipers. The discs are internally vented much like cast-iron ones, and cross-drilled.
Adjustment mechanism
In automotive applications, the piston seal has a square cross-section, also known as a square-cut seal.
As the piston moves in and out, the seal drags and stretches on the piston, causing the seal to twist. The seal distorts approximately 1/10 of a millimeter. The piston is allowed to move out freely, but the slight amount of drag caused by the seal stops the piston from fully retracting to its previous position when the brakes are released, and so takes up the slack caused by the wear of the brake pads, eliminating the need for return springs.
In some rear disc calipers, the parking brake activates a mechanism inside the caliper that performs some of the same functions.
Disc damage modes
Discs are usually damaged in one of four ways: scarring, cracking, warping, or excessive rusting. Service shops will sometimes respond to any disc problem by changing out the discs entirely, This is done mainly where the cost of a new disc may be lower than the cost of labor to resurface the old disc. Mechanically this is unnecessary unless the discs have reached the manufacturer's minimum recommended thickness, which would make it unsafe to use them, or vane rusting is severe (ventilated discs only). Most leading vehicle manufacturers recommend brake disc skimming (US: turning) as a solution for lateral run-out, vibration issues, and brake noises. The machining process is performed in a brake lathe, which removes a very thin layer off the disc surface to clean off minor damage and restore uniform thickness. Machining the disc as necessary will maximize the mileage out of the current discs on the vehicle.
Run-out
Run-out is measured using a dial indicator on a fixed rigid base, with the tip perpendicular to the brake disc's face. It is typically measured about from the outside diameter of the disc. The disc is spun. The difference between the minimum and maximum value on the dial is called lateral run-out. Typical hub/disc assembly run-out specifications for passenger vehicles are around . Runout can be caused either by deformation of the disc itself or by runout in the underlying wheel hub face or by contamination between the disc surface and the underlying hub mounting surface. Determining the root cause of the indicator displacement (lateral runout) requires the disassembly of the disc from the hub. Disc face runout due to hub face runout or contamination will typically have a period of 1 minimum and 1 maximum per revolution of the brake disc.
Discs can be machined to eliminate thickness variation and lateral run-out. Machining can be done in situ (on-car) or off-car (bench lathe). Both methods will eliminate the thickness variation. Machining on-car with the proper equipment can also eliminate lateral run-out due to hub-face non-perpendicularity.
Incorrect fitting can distort (warp) discs. The disc's retaining bolts (or the wheel/lug nuts, if the disc is sandwiched in place by the wheel) must be tightened progressively and evenly. The use of air tools to fasten lug nuts can be bad practice unless a torque wrench is used for final tightening. The vehicle manual will indicate the proper pattern for tightening as well as a torque rating for the bolts. Lug nuts should never be tightened in a circle. Some vehicles are sensitive to the force the bolts apply and tightening should be done with a torque wrench.
Often uneven pad transfer is confused for disc warping. The majority of brake discs diagnosed as "warped" are the result of uneven transfer of pad material. Uneven pad transfer can lead to thickness variation of the disc. When the thicker section of the disc passes between the pads, the pads will move apart and the brake pedal will raise slightly; this is pedal pulsation. The thickness variation can be felt by the driver when it is approximately or greater (on automobile discs).
Thickness variation has many causes, but three primary mechanisms contribute to the propagation of disc thickness variations. The first is the improper selection of brake pads. Pads that are effective at low temperatures, such as when braking for the first time in cold weather, often are made of materials that decompose unevenly at higher temperatures. This uneven decomposition results in the uneven deposition of material onto the brake disc. Another cause of uneven material transfer is the improper break-in of a pad/disc combination. For proper break-in, the disc surface should be refreshed (either by machining the contact surface or by replacing the disc) every time the pads are changed. Once this is done, the brakes are heavily applied multiple times in succession. This creates a smooth, even interface between the pad and the disc. When this is not done properly the brake pads will see an uneven distribution of stress and heat, resulting in an uneven, seemingly random, deposition of pad material. The third primary mechanism of uneven pad material transfer is "pad imprinting." This occurs when the brake pads are heated to the point that the material begins to break down and transfer to the disc. In a properly broken-in brake system (with properly selected pads), this transfer is natural and is a major contributor to the braking force generated by the brake pads. However, if the vehicle comes to a stop and the driver continues to apply the brakes, as is customary in cars with an automatic transmission, the pads will deposit a layer of material in the shape of the brake pad. This small thickness variation can begin the cycle of uneven pad transfer.
Once the disc has some level of variation in thickness, uneven pad deposition can accelerate, sometimes resulting in changes to the crystal structure of the metal that composes the disc. As the brakes are applied, the pads slide over the varying disc surface. As the pads pass by the thicker section of the disc, they are forced outwards. The foot of the driver applied to the brake pedal naturally resists this change, and thus more force is applied to the pads. The result is that the thicker sections see higher levels of stress. This causes uneven heating of the surface of the disc, which causes two major issues. As the brake disc heats unevenly it also expands unevenly. The thicker sections of the disc expand more than the thinner sections due to seeing more heat, and thus the difference in thickness is magnified. Also, the uneven distribution of heat results in the further uneven transfer of pad material. The result is that the thicker-hotter sections receive even more pad material than the thinner-cooler sections, contributing to a further increase in the variation in the disc's thickness. In extreme situations, this uneven heating can cause the crystal structure of the disc material to change. When the hotter sections of the discs reach extremely high temperatures ( ), the metal can undergo a phase transformation and the carbon which is dissolved in the steel can precipitate out to form carbon-heavy carbide regions known as cementite. This iron carbide is very different from the cast iron the rest of the disc is composed of. It is extremely hard, brittle, and does not absorb heat well. After cementite is formed, the integrity of the disc is compromised. Even if the disc surface is machined, the cementite within the disc will not wear or absorb heat at the same rate as the cast iron surrounding it, causing the uneven thickness and heating characteristics of the disc to return.
Scarring
Scarring (US: Scoring) can occur if brake pads are not changed promptly when they reach the end of their service life and are considered worn out. Once enough of the friction material has worn away, the pad's steel backing plate (for glued pads) or the pad retainer rivets (for riveted pads) will bear upon the disc's wear surface, reducing braking power and making scratches on the disc. Generally, a moderately scarred or scored disc, which operated satisfactorily with existing brake pads, will be equally usable with new pads. If the scarring is deeper but not excessive, it can be repaired by machining off a layer of the disc's surface. This can only be done a limited number of times as the disc has a minimum rated safe thickness. The minimum thickness value is typically cast into the disc during manufacturing on the hub or the edge of the disc. In Pennsylvania, which has one of the most rigorous auto safety inspection programs in North America, an automotive disc cannot pass a safety inspection if any scoring is deeper than , and must be replaced if machining will reduce the disc below its minimum safe thickness.
To prevent scarring, it is prudent to periodically inspect the brake pads for wear. A tire rotation is a logical time for inspection, since rotation must be performed regularly based on vehicle operation time and all wheels must be removed, allowing ready visual access to the brake pads. Some types of alloy wheels and brake arrangements will provide enough open space to view the pads without removing the wheel. When practical, pads that are near the wear-out point should be replaced immediately, as complete wear-out leads to scarring damage and unsafe braking. Many disc brake pads will include some sort of soft steel spring or drag tab as part of the pad assembly, which drags on the disc when the pad is nearly worn out. This produces a moderately loud squealing noise, alerting the driver that service is required. This will not normally scar the disc if the brakes are serviced promptly. A set of pads can be considered for replacement if the thickness of the pad material is the same or less than the thickness of the backing steel. In Pennsylvania, the standard is for riveted pads and 2/32" for bonded pads.
Cracking
Cracking is limited mostly to drilled discs, which may develop small cracks around the edges of holes drilled near the edge of the disc due to the disc's uneven rate of expansion in severe-duty environments. Manufacturers that use drilled discs as OEM typically do so for two reasons: appearance, if they determine that the average owner of the vehicle model will prefer the look while not overly stressing the hardware; or as a function of reducing the unsprung weight of the brake assembly, with the engineering assumption that enough brake disc mass remains to absorb racing temperatures and stresses. A brake disc is a heat sink, but the loss of heat sink mass may be balanced by increased surface area to radiate away heat. Small hairline cracks may appear in any cross-drilled metal disc as a normal wear mechanism, but in severe cases, the disc will fail catastrophically. No repair is possible for the cracks, and if the cracking becomes severe, the disc must be replaced. These cracks occur due to the phenomenon of low cycle fatigue as a result of repeated hard braking.
Rusting
The discs are commonly made from cast iron and a certain amount of surface rust is normal. The disc contact area for the brake pads will be kept clean by regular use, but a vehicle that is stored for an extended period can develop significant rust in the contact area that may reduce braking power for a time until the rusted layer is worn off again. Rusting can also lead to disc warping when brakes are re-activated after storage because of differential heating between unrusted areas left covered by pads and rust around the majority of the disc area surface. Over time, vented brake discs may develop severe rust corrosion inside the ventilation slots, compromising the strength of the structure and needing replacement.
Calipers
The brake caliper is the assembly that houses the brake pads and pistons. The pistons are usually made of plastic, aluminium or chrome-plated steel.
Calipers are of two types, floating or fixed. A fixed caliper does not move relative to the disc and is thus less tolerant of disc imperfections. It uses one or more pairs of opposing pistons to clamp from each side of the disc and is more complex and expensive than a floating caliper.
A floating caliper (also called a "sliding caliper") moves side to side to the disc, along a line parallel to the axis of rotation of the disc; a piston on one side of the disc pushes the inner brake pad until it makes contact with the braking surface, then pulls the caliper body with the outer brake pad so the pressure is applied to both sides of the disc. Floating caliper (single piston) designs are subject to sticking failure, caused by dirt or corrosion entering at least one mounting mechanism and stopping its normal movement. This can lead to the caliper's pads rubbing on the disc when the brake is not engaged or engaging it at an angle. Sticking can result from infrequent vehicle use, failure of a seal or rubber protection boot allowing debris entry, dry-out of the grease in the mounting mechanism, and subsequent moisture incursion leading to corrosion, or some combination of these factors. Consequences may include reduced fuel efficiency, extreme heating of the disc, or excessive wear on the affected pad. A sticking front caliper may also cause steering vibration.
Another type of floating caliper is a swinging caliper. Instead of a pair of horizontal bolts that allow the caliper to move straight in and out respective to the car body, a swinging caliper utilizes a single, vertical pivot bolt located somewhere behind the axle centerline. When the driver presses the brakes, the brake piston pushes on the inside piston and rotates the whole caliper inward, when viewed from the top. Because the swinging caliper's piston angle changes relative to the disc, this design uses wedge-shaped pads that are narrower in the rear on the outside and narrower in the front on the inside.
Various types of brake calipers are also used on bicycle rim brakes.
Pistons and cylinders
The most common caliper design uses a single hydraulically actuated piston within a cylinder, although high-performance brakes use as many as twelve. Modern cars use different hydraulic circuits to actuate the brakes on each set of wheels as a safety measure. The hydraulic design also helps multiply braking force. The number of pistons in a caliper is often referred to as the number of 'pots', so if a vehicle has 'six pot' calipers it means that each caliper houses six pistons.
Brake failure can result from the failure of the piston to retract, which is usually a consequence of not operating the vehicle during prolonged storage outdoors in adverse conditions. On high-mileage vehicles, the piston seals may leak, which must be promptly corrected.
Brake pads
Brake pads are designed for high friction with brake pad material embedded in the disc in the process of bedding while wearing evenly. Friction can be divided into two parts. They are: adhesive and abrasive.
Depending on the properties of the material of both the pad and the disc and the configuration and the usage, pad and disc wear rates will vary considerably. The properties that determine material wear involve trade-offs between performance and longevity.
The brake pads must usually be replaced regularly (depending on pad material and driving style), and some are equipped with a mechanism that alerts drivers that replacement is needed, such as a thin piece of soft metal that rubs against the disc when the pads are too thin causing the brakes to squeal, a soft metal tab embedded in the pad material that closes an electric circuit and lights a warning light when the brake pad gets thin, or an electronic sensor.
Generally, road-going vehicles have two brake pads per caliper, while up to six are installed on each racing caliper, with varying frictional properties in a staggered pattern for optimum performance.
Early brake pads (and linings) contained asbestos, producing dust that should not be inhaled. Although newer pads can be made of ceramics, Kevlar, and other plastics, inhalation of brake dust should still be avoided regardless of material.
Common problems
Squeal
Sometimes a loud noise or high-pitched squeal occurs when the brakes are applied. It mostly occurs on old cars and those who were produced or acquired some time ago. Most brake squeal is produced by vibration (resonance instability) of the brake components, especially the pads and discs (known as force-coupled excitation). This type of squeal should not negatively affect brake-stopping performance. Techniques include adding chamfer pads to the contact points between the caliper pistons and the pads, the bonding insulators (damping material) to the pad backplate, the brake shims between the brake pad and pistons, etc. All should be coated with an extremely high temperature, high solids lubricant to help reduce squeal. This allows the metal-to-metal parts to move independently of each other and thereby eliminate the buildup of energy that can create a frequency that is heard as brake squeal, groan, or growl. It is inherent that some pads are going to squeal more given the type of the pad and its usage case. Pads typically rated to withstand very high temperatures for extended periods tend to produce high amounts of friction leading to more noise during brake application.
Cold weather combined with high early-morning humidity (dew) often worsens brake squeal. However, the squeal generally stops when the lining reaches regular operating temperatures. This more strongly affects pads meant to be used at higher temperatures. Dust on the brakes may also cause squeal and commercial brake cleaning products are designed to remove dirt and other contaminants. Pads without a proper amount of transfer material could also squeal, this can be remedied by bedding or re-bedding the brake pads to brake discs.
Some lining wear indicators, located either as a semi-metallic layer within the brake pad material or with an external "sensor", are also designed to squeal when the lining is due for replacement. The typical external sensor is fundamentally different from the noises described above (when the brakes are applied) because the wear sensor noise typically occurs when the brakes are not used. The wear sensor may only create a squeal under braking when it first begins to indicate wear but is still a fundamentally different sound and pitch.
Judder or shimmy
A brake judder is usually perceived by the driver as minor to severe vibrations transferred through the chassis during braking.
The judder phenomenon can be classified into two distinct subgroups: hot (or thermal), or cold judder.
Hot judder is usually produced as a result of longer, more moderate braking from high speed where the vehicle does not come to a complete stop. It commonly occurs when a motorist decelerates from speeds of around to about , which results in severe vibrations being transmitted to the driver. These vibrations are the result of uneven thermal distributions, or hot spots. Hot spots are classified as concentrated thermal regions that alternate between both sides of a disc that distort it in such a way that produces a sinusoidal waviness around its edges. Once the brake pads (friction material/brake lining) come in contact with the sinusoidal surface during braking, severe vibrations are induced, and can produce hazardous conditions for the person driving the vehicle.
Cold judder, on the other hand, is the result of uneven disc wear patterns or disc thickness variation (DTV). These variations in the disc surface are usually the result of extensive vehicle road usage. DTV is usually attributed to the following causes: waviness and roughness of disc surface, misalignment of axis run-out, elastic deflection, wear and friction material transfers.
Either type could potentially be fixed by ensuring a clean mounting surface on either side of the brake disc between the wheel hub and brake disc hub before usage and paying attention to imprinting after extended usage by leaving the brake pedal heavily depressed at the end of heavy usage. Sometimes a bed in procedure can clean and minimize DTV and lay a new even transfer layer between the pad and brake disc. However, it will not eliminate hot spots or excessive run-out.
Dust
When braking force is applied, the act of abrasive friction between the brake pad and the disc wears both the disc and pad away. The brake dust that is seen deposited on wheels, calipers, and other braking system components consists mostly of the disc material. Brake dust can damage the finish of most wheels if not washed off. Generally, a brake pad that aggressively abrades more disc material away, such as metallic pads, will create more brake dust. Some higher-performing pads for track use or towing use may wear away much quicker than a typical pad, thus causing more dust due to the increased brake disc and brake pad wear.
Brake fade
Brake fade is a phenomenon that decreases braking efficiency. It causes the braking power to reduce and you feel that the brakes are not being applied with the force they were being applied at the time of starting. This occurs due to the heating of brake pads. The heated brake pads emit some gaseous substances which cover the area between the disc and the brake pads. These gases disturb the contact between the brake pads and the disc and hence decrease the braking effectiveness.
Patents
| Technology | Motorized road transport | null |
59611 | https://en.wikipedia.org/wiki/Ionization | Ionization | Ionization (or ionisation specifically in Britain, Ireland, Australia and New Zealand) is the process by which an atom or a molecule acquires a negative or positive charge by gaining or losing electrons, often in conjunction with other chemical changes. The resulting electrically charged atom or molecule is called an ion. Ionization can result from the loss of an electron after collisions with subatomic particles, collisions with other atoms, molecules, electrons, positrons, protons, antiprotons and ions, or through the interaction with electromagnetic radiation. Heterolytic bond cleavage and heterolytic substitution reactions can result in the formation of ion pairs. Ionization can occur through radioactive decay by the internal conversion process, in which an excited nucleus transfers its energy to one of the inner-shell electrons causing it to be ejected.
Uses
Everyday examples of gas ionization occur within a fluorescent lamp or other electrical discharge lamps. It is also used in radiation detectors such as the Geiger-Müller counter or the ionization chamber. The ionization process is widely used in a variety of equipment in fundamental science (e.g., mass spectrometry) and in medical treatment (e.g., radiation therapy). It is also widely used for air purification, though studies have shown harmful effects of this application.
Production of ions
Negatively charged ions are produced when a free electron collides with an atom and is subsequently trapped inside the electric potential barrier, releasing any excess energy. The process is known as electron capture ionization.
Positively charged ions are produced by transferring an amount of energy to a bound electron in a collision with charged particles (e.g. ions, electrons or positrons) or with photons. The threshold amount of the required energy is known as ionization potential. The study of such collisions is of fundamental importance with regard to the few-body problem, which is one of the major unsolved problems in physics. Kinematically complete experiments, i.e. experiments in which the complete momentum vector of all collision fragments (the scattered projectile, the recoiling target-ion, and the ejected electron) are determined, have contributed to major advances in the theoretical understanding of the few-body problem in recent years.
Adiabatic ionization
Adiabatic ionization is a form of ionization in which an electron is removed from or added to an atom or molecule in its lowest energy state to form an ion in its lowest energy state.
The Townsend discharge is a good example of the creation of positive ions and free electrons due to ion impact. It is a cascade reaction involving electrons in a region with a sufficiently high electric field in a gaseous medium that can be ionized, such as air. Following an original ionization event, due to such as ionizing radiation, the positive ion drifts towards the cathode, while the free electron drifts towards the anode of the device. If the electric field is strong enough, the free electron gains sufficient energy to liberate a further electron when it next collides with another molecule. The two free electrons then travel towards the anode and gain sufficient energy from the electric field to cause impact ionization when the next collisions occur; and so on. This is effectively a chain reaction of electron generation, and is dependent on the free electrons gaining sufficient energy between collisions to sustain the avalanche.
Ionization efficiency is the ratio of the number of ions formed to the number of electrons or photons used.
Ionization energy of atoms
The trend in the ionization energy of atoms is often used to demonstrate the periodic behavior of atoms with respect to the atomic number, as summarized by ordering atoms in Mendeleev's table. This is a valuable tool for establishing and understanding the ordering of electrons in atomic orbitals without going into the details of wave functions or the ionization process. An example is presented in the figure to the right. The periodic abrupt decrease in ionization potential after rare gas atoms, for instance, indicates the emergence of a new shell in alkali metals. In addition, the local maximums in the ionization energy plot, moving from left to right in a row, are indicative of s, p, d, and f sub-shells.
Semi-classical description of ionization
Classical physics and the Bohr model of the atom can qualitatively explain photoionization and collision-mediated ionization. In these cases, during the ionization process, the energy of the electron exceeds the energy difference of the potential barrier it is trying to pass. The classical description, however, cannot describe tunnel ionization since the process involves the passage of electron through a classically forbidden potential barrier.
Quantum mechanical description of ionization
The interaction of atoms and molecules with sufficiently strong laser pulses or with other charged particles leads to the ionization to singly or multiply charged ions. The ionization rate, i.e. the ionization probability in unit time, can be calculated using quantum mechanics. (There are classical methods available also, like the Classical Trajectory Monte Carlo Method (CTMC), but it is not overall accepted and often criticized by the community.) There are two quantum mechanical methods exist, perturbative and non-perturbative methods like time-dependent coupled-channel or time independent close coupling methods where the wave function is expanded in a finite basis set. There are numerous options available e.g. B-splines, generalized Sturmians or Coulomb wave packets. Another non-perturbative method is to solve the corresponding Schrödinger equation fully numerically on a lattice.
In general, the analytic solutions are not available, and the approximations required for manageable numerical calculations do not provide accurate enough results. However, when the laser intensity is sufficiently high, the detailed structure of the atom or molecule can be ignored and analytic solution for the ionization rate is possible.
Tunnel ionization
Tunnel ionization is ionization due to quantum tunneling. In classical ionization, an electron must have enough energy to make it over the potential barrier, but quantum tunneling allows the electron simply to go through the potential barrier instead of going all the way over it because of the wave nature of the electron. The probability of an electron's tunneling through the barrier drops off exponentially with the width of the potential barrier. Therefore, an electron with a higher energy can make it further up the potential barrier, leaving a much thinner barrier to tunnel through and thus a greater chance to do so. In practice, tunnel ionization is observable when the atom or molecule is interacting with near-infrared strong laser pulses. This process can be understood as a process by which a bounded electron, through the absorption of more than one photon from the laser field, is ionized. This picture is generally known as multiphoton ionization (MPI).
Keldysh modeled the MPI process as a transition of the electron from the ground state of the atom to the Volkov states. In this model the perturbation of the ground state by the laser field is neglected and the details of atomic structure in determining the ionization probability are not taken into account. The major difficulty with Keldysh's model was its neglect of the effects of Coulomb interaction on the final state of the electron. As it is observed from figure, the Coulomb field is not very small in magnitude compared to the potential of the laser at larger distances from the nucleus. This is in contrast to the approximation made by neglecting the potential of the laser at regions near the nucleus. Perelomov et al. included the Coulomb interaction at larger internuclear distances. Their model (which we call the PPT model) was derived for short range potential and includes the effect of the long range Coulomb interaction through the first order correction in the quasi-classical action. Larochelle et al. have compared the theoretically predicted ion versus intensity curves of rare gas atoms interacting with a Ti:Sapphire laser with experimental measurement. They have shown that the total ionization rate predicted by the PPT model fit very well the experimental ion yields for all rare gases in the intermediate regime of the Keldysh parameter.
The rate of MPI on atom with an ionization potential in a linearly polarized laser with frequency is given by
where
is the Keldysh parameter,
,
is the peak electric field of the laser and
.
The coefficients , and are given by
The coefficient is given by
where
Quasi-static tunnel ionization
The quasi-static tunneling (QST) is the ionization whose rate can be satisfactorily predicted by the ADK model, i.e. the limit of the PPT model when approaches zero. The rate of QST is given by
As compared to the absence of summation over n, which represent different above threshold ionization (ATI) peaks, is remarkable.
Strong field approximation for the ionization rate
The calculations of PPT are done in the E-gauge, meaning that the laser field is taken as electromagnetic waves. The ionization rate can also be calculated in A-gauge, which emphasizes the particle nature of light (absorbing multiple photons during ionization). This approach was adopted by Krainov model based on the earlier works of Faisal and Reiss. The resulting rate is given by
where:
with being the ponderomotive energy,
is the minimum number of photons necessary to ionize the atom,
is the double Bessel function,
with the angle between the momentum of the electron, p, and the electric field of the laser, F,
FT is the three-dimensional Fourier transform, and
incorporates the Coulomb correction in the SFA model.
Population trapping
In calculating the rate of MPI of atoms only transitions to the continuum states are considered. Such an approximation is acceptable as long as there is no multiphoton resonance between the ground state and some excited states. However, in real situation of interaction with pulsed lasers, during the evolution of laser intensity, due to different Stark shift of the ground and excited states there is a possibility that some excited state go into multiphoton resonance with the ground state. Within the dressed atom picture, the ground state dressed by photons and the resonant state undergo an avoided crossing at the resonance intensity . The minimum distance, , at the avoided crossing is proportional to the generalized Rabi frequency, coupling the two states. According to Story et al., the probability of remaining in the ground state, , is given by
where is the time-dependent energy difference between the two dressed states. In interaction with a short pulse, if the dynamic resonance is reached in the rising or the falling part of the pulse, the population practically remains in the ground state and the effect of multiphoton resonances may be neglected. However, if the states go onto resonance at the peak of the pulse, where , then the excited state is populated. After being populated, since the ionization potential of the excited state is small, it is expected that the electron will be instantly ionized.
In 1992, de Boer and Muller showed that Xe atoms subjected to short laser pulses could survive in the highly excited states 4f, 5f, and 6f. These states were believed to have been excited by the dynamic Stark shift of the levels into multiphoton resonance with the field during the rising part of the laser pulse. Subsequent evolution of the laser pulse did not completely ionize these states, leaving behind some highly excited atoms. We shall refer to this phenomenon as "population trapping".
We mention the theoretical calculation that incomplete ionization occurs whenever there is parallel resonant excitation into a common level with ionization loss. We consider a state such as 6f of Xe which consists of 7 quasi-degnerate levels in the range of the laser bandwidth. These levels along with the continuum constitute a lambda system. The mechanism of the lambda type trapping is schematically presented in figure. At the rising part of the pulse (a) the excited state (with two degenerate levels 1 and 2) are not in multiphoton resonance with the ground state. The electron is ionized through multiphoton coupling with the continuum. As the intensity of the pulse is increased the excited state and the continuum are shifted in energy due to the Stark shift. At the peak of the pulse (b) the excited states go into multiphoton resonance with the ground state. As the intensity starts to decrease (c), the two state are coupled through continuum and the population is trapped in a coherent superposition of the two states. Under subsequent action of the same pulse, due to interference in the transition amplitudes of the lambda system, the field cannot ionize the population completely and a fraction of the population will be trapped in a coherent superposition of the quasi degenerate levels. According to this explanation the states with higher angular momentum – with more sublevels – would have a higher probability of trapping the population. In general the strength of the trapping will be determined by the strength of the two photon coupling between the quasi-degenerate levels via the continuum. In 1996, using a very stable laser and by minimizing the masking effects of the focal region expansion with increasing intensity, Talebpour et al. observed structures on the curves of singly charged ions of Xe, Kr and Ar. These structures were attributed to electron trapping in the strong laser field. A more unambiguous demonstration of population trapping has been reported by T. Morishita and C. D. Lin.
Non-sequential multiple ionization
The phenomenon of non-sequential ionization (NSI) of atoms exposed to intense laser fields has been a subject of many theoretical and experimental studies since 1983. The pioneering work began with the observation of a "knee" structure on the Xe2+ ion signal versus intensity curve by L’Huillier et al. From the experimental point of view, the NS double ionization refers to processes which somehow enhance the rate of production of doubly charged ions by a huge factor at intensities below the saturation intensity of the singly charged ion. Many, on the other hand, prefer to define the NSI as a process by which two electrons are ionized nearly simultaneously. This definition implies that apart from the sequential channel there is another channel which is the main contribution to the production of doubly charged ions at lower intensities. The first observation of triple NSI in argon interacting with a 1 μm laser was reported by Augst et al. Later, systematically studying the NSI of all rare gas atoms, the quadruple NSI of Xe was observed. The most important conclusion of this study was the observation of the following relation between the rate of NSI to any charge state and the rate of tunnel ionization (predicted by the ADK formula) to the previous charge states;
where is the rate of quasi-static tunneling to i'th charge state and are some constants depending on the wavelength of the laser (but not on the pulse duration).
Two models have been proposed to explain the non-sequential ionization; the shake-off model and electron re-scattering model. The shake-off (SO) model, first proposed by Fittinghoff et al., is adopted from the field of ionization of atoms by X rays and electron projectiles where the SO process is one of the major mechanisms responsible for the multiple ionization of atoms. The SO model describes the NSI process as a mechanism where one electron is ionized by the laser field and the departure of this electron is so rapid that the remaining electrons do not have enough time to adjust themselves to the new energy states. Therefore, there is a certain probability that, after the ionization of the first electron, a second electron is excited to states with higher energy (shake-up) or even ionized (shake-off). We should mention that, until now, there has been no quantitative calculation based on the SO model, and the model is still qualitative.
The electron rescattering model was independently developed by Kuchiev, Schafer et al, Corkum, Becker and Faisal and Faisal and Becker. The principal features of the model can be understood easily from Corkum's version. Corkum's model describes the NS ionization as a process whereby an electron is tunnel ionized. The electron then interacts with the laser field where it is accelerated away from the nuclear core. If the electron has been ionized at an appropriate phase of the field, it will pass by the position of the remaining ion half a cycle later, where it can free an additional electron by electron impact. Only half of the time the electron is released with the appropriate phase and the other half it never return to the nuclear core. The maximum kinetic energy that the returning electron can have is 3.17 times the ponderomotive potential () of the laser. Corkum's model places a cut-off limit on the minimum intensity ( is proportional to intensity) where ionization due to re-scattering can occur.
The re-scattering model in Kuchiev's version (Kuchiev's model) is quantum mechanical. The basic idea of the model is illustrated by Feynman diagrams in figure a. First both electrons are in the ground state of an atom. The lines marked a and b describe the corresponding atomic states. Then the electron a is ionized. The beginning of the ionization process is shown by the intersection with a sloped dashed line. where the MPI occurs. The propagation of the ionized electron in the laser field, during which it absorbs other photons (ATI), is shown by the full thick line. The collision of this electron with the parent atomic ion is shown by a vertical dotted line representing the Coulomb interaction between the electrons. The state marked with c describes the ion excitation to a discrete or continuum state. Figure b describes the exchange process. Kuchiev's model, contrary to Corkum's model, does not predict any threshold intensity for the occurrence of NS ionization.
Kuchiev did not include the Coulomb effects on the dynamics of the ionized electron. This resulted in the underestimation of the double ionization rate by a huge factor. Obviously, in the approach of Becker and Faisal (which is equivalent to Kuchiev's model in spirit), this drawback does not exist. In fact, their model is more exact and does not suffer from the large number of approximations made by Kuchiev. Their calculation results perfectly fit with the experimental results of Walker et al. Becker and Faisal have been able to fit the experimental results on the multiple NSI of rare gas atoms using their model. As a result, the electron re-scattering can be taken as the main mechanism for the occurrence of the NSI process.
Multiphoton ionization of inner-valence electrons and fragmentation of polyatomic molecules
The ionization of inner valence electrons are responsible for the fragmentation of polyatomic molecules in strong laser fields. According to a qualitative model the dissociation of the molecules occurs through a three-step mechanism:
MPI of electrons from the inner orbitals of the molecule which results in a molecular ion in ro-vibrational levels of an excited electronic state;
Rapid radiationless transition to the high-lying ro-vibrational levels of a lower electronic state; and
Subsequent dissociation of the ion to different fragments through various fragmentation channels.
The short pulse induced molecular fragmentation may be used as an ion source for high performance mass spectroscopy. The selectivity provided by a short pulse based source is superior to that expected when using the conventional electron ionization based sources, in particular when the identification of optical isomers is required.
Kramers–Henneberger frame
The Kramers–Henneberger(KF) frame is the non-inertial frame moving with the free electron under the influence of the harmonic laser pulse, obtained by applying a translation to the laboratory frame equal to the quiver motion of a classical electron in the laboratory frame. In other words, in the Kramers–Henneberger frame the classical electron is at rest. Starting in the lab frame (velocity gauge), we may describe the electron with the Hamiltonian:
In the dipole approximation, the quiver motion of a classical electron in the laboratory frame for an arbitrary field can be obtained from the vector potential of the electromagnetic field:
where for a monochromatic plane wave.
By applying a transformation to the laboratory frame equal to the quiver motion one moves to the ‘oscillating’ or ‘Kramers–Henneberger’ frame, in which the classical electron is at rest. By a phase factor transformation for convenience one obtains the ‘space-translated’ Hamiltonian, which is unitarily equivalent to the lab-frame Hamiltonian, which contains the original potential centered on the oscillating point :
The utility of the KH frame lies in the fact that in this frame the laser-atom interaction can be reduced to the form of an oscillating potential energy, where the natural parameters describing the electron dynamics are and (sometimes called the “excursion amplitude’, obtained from ).
From here one can apply Floquet theory to calculate quasi-stationary solutions of the TDSE. In high frequency Floquet theory, to lowest order in the system reduces to the so-called ‘structure equation’, which has the form of a typical energy-eigenvalue Schrödinger equation containing the ‘dressed potential’ (the cycle-average of the oscillating potential). The interpretation of the presence of is as follows: in the oscillating frame, the nucleus has an oscillatory motion of trajectory and can be seen as the potential of the smeared out nuclear charge along its trajectory.
The KH frame is thus employed in theoretical studies of strong-field ionization and atomic stabilization (a predicted phenomenon in which the ionization probability of an atom in a high-intensity, high-frequency field actually decreases for intensities above a certain threshold) in conjunction with high-frequency Floquet theory.
The KF frame was successfully applied for different problems as well e.g. for higher-hamonic generation from a metal surface in a powerful laser field
Dissociation – distinction
A substance may dissociate without necessarily producing ions. As an example, the molecules of table sugar dissociate in water (sugar is dissolved) but exist as intact neutral entities. Another subtle event is the dissociation of sodium chloride (table salt) into sodium and chlorine ions. Although it may seem as a case of ionization, in reality the ions already exist within the crystal lattice. When salt is dissociated, its constituent ions are simply surrounded by water molecules and their effects are visible (e.g. the solution becomes electrolytic). However, no transfer or displacement of electrons occurs.
| Physical sciences | Phase transitions | null |
59613 | https://en.wikipedia.org/wiki/Ionization%20energy | Ionization energy | In physics and chemistry, ionization energy (IE) is the minimum energy required to remove the most loosely bound electron of an isolated gaseous atom, positive ion, or molecule. The first ionization energy is quantitatively expressed as
X(g) + energy ⟶ X+(g) + e−
where X is any atom or molecule, X+ is the resultant ion when the original atom was stripped of a single electron, and e− is the removed electron. Ionization energy is positive for neutral atoms, meaning that the ionization is an endothermic process. Roughly speaking, the closer the outermost electrons are to the nucleus of the atom, the higher the atom's ionization energy.
In physics, ionization energy is usually expressed in electronvolts (eV) or joules (J). In chemistry, it is expressed as the energy to ionize a mole of atoms or molecules, usually as kilojoules per mole (kJ/mol) or kilocalories per mole (kcal/mol).
Comparison of ionization energies of atoms in the periodic table reveals two periodic trends which follow the rules of Coulombic attraction:
Ionization energy generally increases from left to right within a given period (that is, row).
Ionization energy generally decreases from top to bottom in a given group (that is, column).
The latter trend results from the outer electron shell being progressively farther from the nucleus, with the addition of one inner shell per row as one moves down the column.
The nth ionization energy refers to the amount of energy required to remove the most loosely bound electron from the species having a positive charge of (n − 1). For example, the first three ionization energies are defined as follows:
1st ionization energy is the energy that enables the reaction X ⟶ X+ + e−
2nd ionization energy is the energy that enables the reaction X+ ⟶ X2+ + e−
3rd ionization energy is the energy that enables the reaction X2+ ⟶ X3+ + e−
The most notable influences that determine ionization energy include:
Electron configuration: This accounts for most elements' IE, as all of their chemical and physical characteristics can be ascertained just by determining their respective electron configuration.
Nuclear charge: If the nuclear charge (atomic number) is greater, the electrons are held more tightly by the nucleus and hence the ionization energy will be greater (leading to the mentioned trend 1 within a given period).
Number of electron shells: If the size of the atom is greater due to the presence of more shells, the electrons are held less tightly by the nucleus and the ionization energy will be smaller.
Effective nuclear charge (Zeff): If the magnitude of electron shielding and penetration are greater, the electrons are held less tightly by the nucleus, the Zeff of the electron and the ionization energy is smaller.
Stability: An atom having a more stable electronic configuration has a reduced tendency to lose electrons and consequently has a higher ionization energy.
Minor influences include:
Relativistic effects: Heavier elements (especially those whose atomic number is greater than about 70) are affected by these as their electrons are approaching the speed of light. They therefore have smaller atomic radii and higher ionization energies.
Lanthanide and actinide contraction (and scandide contraction): The shrinking of the elements affects the ionization energy, as the net charge of the nucleus is more strongly felt.
Electron pairing energies: Half-filled subshells usually result in higher ionization energies.
The term ionization potential is an older and obsolete term for ionization energy, because the oldest method of measuring ionization energy was based on ionizing a sample and accelerating the electron removed using an electrostatic potential.
Determination of ionization energies
The ionization energy of atoms, denoted Ei, is measured by finding the minimal energy of light quanta (photons) or electrons accelerated to a known energy that will kick out the least bound atomic electrons. The measurement is performed in the gas phase on single atoms. While only noble gases occur as monatomic gases, other gases can be split into single atoms. Also, many solid elements can be heated and vaporized into single atoms. Monatomic vapor is contained in a previously evacuated tube that has two parallel electrodes connected to a voltage source. The ionizing excitation is introduced through the walls of the tube or produced within.
When ultraviolet light is used, the wavelength is swept down the ultraviolet range. At a certain wavelength (λ) and frequency of light (ν=c/λ, where c is the speed of light), the light quanta, whose energy is proportional to the frequency, will have energy high enough to dislodge the least bound electrons. These electrons will be attracted to the positive electrode, and the positive ions remaining after the photoionization will get attracted to the negatively charged electrode. These electrons and ions will establish a current through the tube. The ionization energy will be the energy of photons hνi (h is the Planck constant) that caused a steep rise in the current: Ei = hνi.
When high-velocity electrons are used to ionize the atoms, they are produced by an electron gun inside a similar evacuated tube. The energy of the electron beam can be controlled by the acceleration voltages. The energy of these electrons that gives rise to a sharp onset of the current of ions and freed electrons through the tube will match the ionization energy of the atoms.
Atoms: values and trends
Generally, the (N+1)th ionization energy of a particular element is larger than the Nth ionization energy (it may also be noted that the ionization energy of an anion is generally less than that of cations and neutral atom for the same element). When the next ionization energy involves removing an electron from the same electron shell, the increase in ionization energy is primarily due to the increased net charge of the ion from which the electron is being removed. Electrons removed from more highly charged ions experience greater forces of electrostatic attraction; thus, their removal requires more energy. In addition, when the next ionization energy involves removing an electron from a lower electron shell, the greatly decreased distance between the nucleus and the electron also increases both the electrostatic force and the distance over which that force must be overcome to remove the electron. Both of these factors further increase the ionization energy.
Some values for elements of the third period are given in the following table:
Large jumps in the successive molar ionization energies occur when passing noble gas configurations. For example, as can be seen in the table above, the first two molar ionization energies of magnesium (stripping the two 3s electrons from a magnesium atom) are much smaller than the third, which requires stripping off a 2p electron from the neon configuration of Mg2+. That 2p electron is much closer to the nucleus than the 3s electrons removed previously.
Ionization energy is also a periodic trend within the periodic table. Moving left to right within a period, or upward within a group, the first ionization energy generally increases, with exceptions such as aluminium and sulfur in the table above. As the nuclear charge of the nucleus increases across the period, the electrostatic attraction increases between electrons and protons, hence the atomic radius decreases, and the electron cloud comes closer to the nucleus because the electrons, especially the outermost one, are held more tightly by the higher effective nuclear charge.
On moving downward within a given group, the electrons are held in higher-energy shells with higher principal quantum number n, further from the nucleus and therefore are more loosely bound so that the ionization energy decreases. The effective nuclear charge increases only slowly so that its effect is outweighed by the increase in n.
Exceptions in ionization energies
There are exceptions to the general trend of rising ionization energies within a period. For example, the value decreases from beryllium (: 9.3 eV) to boron (: 8.3 eV), and from nitrogen (: 14.5 eV) to oxygen (: 13.6 eV). These dips can be explained in terms of electron configurations.
Boron has its last electron in a 2p orbital, which has its electron density further away from the nucleus on average than the 2s electrons in the same shell. The 2s electrons then shield the 2p electron from the nucleus to some extent, and it is easier to remove the 2p electron from boron than to remove a 2s electron from beryllium, resulting in a lower ionization energy for B.
In oxygen, the last electron shares a doubly occupied p-orbital with an electron of opposing spin. The two electrons in the same orbital are closer together on average than two electrons in different orbitals, so that they shield each other from the nucleus more effectively and it is easier to remove one electron, resulting in a lower ionization energy.
Furthermore, after every noble gas element, the ionization energy drastically drops. This occurs because the outer electron in the alkali metals requires a much lower amount of energy to be removed from the atom than the inner shells. This also gives rise to low electronegativity values for the alkali metals.
The trends and exceptions are summarized in the following subsections:
Ionization energy decreases when
Transitioning to a new period: an alkali metal easily loses one electron to leave an octet or pseudo-noble gas configuration, so those elements have only small values for IE.
Moving from the s-block to the p-block: a p-orbital loses an electron more easily. An example is beryllium to boron, with electron configuration 1s2 2s2 2p1. The 2s electrons shield the higher-energy 2p electron from the nucleus, making it slightly easier to remove. This also happens from magnesium to aluminium.
Occupying a p-subshell with its first electron with spin opposed to the other electrons: such as in nitrogen (: 14.5 eV) to oxygen (: 13.6 eV), as well as phosphorus (: 10.48 eV) to sulfur (: 10.36 eV). The reason for this is because oxygen, sulfur and selenium all have dipping ionization energies because of shielding effects. However, this discontinues starting from tellurium where the shielding is too small to produce a dip.
Moving from the d-block to the p-block: as in the case of zinc (: 9.4 eV) to gallium (: 6.0 eV)
Special case: decrease from lead (: 7.42 eV) to bismuth (: 7.29 eV). This cannot be attributed to size (the difference is minimal: lead has a covalent radius of 146 pm whereas bismuth's is 148 pm). This is due to the spin-orbit splitting of the 6p shell (lead is removing an electron from the stabilised 6p1/2 level, but bismuth is removing one from the destabilised 6p3/2 level). Predicted ionization energies show a much greater decrease from flerovium to moscovium, one row further down the periodic table and with much larger spin-orbit effects.
Special case: decrease from radium (: 5.27 eV) to actinium (: 5.17 eV), which is a switch from an s to a d orbital. However the analogous switch from barium (: 5.2 eV) to lanthanum (: 5.6 eV) does not show a downward change.
Lutetium () and lawrencium () both have ionization energies lower than the previous elements. In both cases the last electron added starts a new subshell: 5d for Lu with electron configuration [Xe] 4f14 5d1 6s2, and 7p for Lr with configuration [Rn] 5f4 7s2 7p1. These dips in ionization energies for lutetium and especially lawrencium show that these elements belong in the d-block, and not lanthanum and actinium.
Ionization energy increases when
Reaching Group 18 noble gas elements: This is due to their complete electron subshells, so that these elements require large amounts of energy to remove one electron.
Group 12: The elements here, zinc (: 9.4 eV), cadmium (: 9.0 eV) and mercury (: 10.4 eV) all record sudden rising IE values in contrast to their preceding elements: copper (: 7.7 eV), silver (: 7.6 eV) and gold (: 9.2 eV), respectively. For mercury, it can be extrapolated that the relativistic stabilization of the 6s electrons increases the ionization energy, in addition to poor shielding by 4f electrons that increases the effective nuclear charge on the outer valence electrons. In addition, the closed-subshells electron configurations: [Ar] 3d10 4s2, [Kr] 4d105s2 and [Xe] 4f14 5d10 6s2 provide increased stability.
Special case: shift from rhodium (: 7.5 eV) to palladium (: 8.3 eV). Unlike other Group 10 elements, palladium has a higher ionization energy than the preceding atom, due to its electron configuration. In contrast to nickel's [Ar] 3d8 4s2, and platinum's [Xe] 4f14 5d9 6s1, palladium's electron configuration is [Kr] 4d10 5s0 (even though the Madelung rule predicts [Kr] 4d8 5s2). Finally, silver's lower IE (: 7.6 eV) further accentuates the high value for palladium; the single added s electron is removed with a lower ionization energy than palladium, which emphasizes palladium's high IE (as shown in the above linear table values for IE)
The IE of gadolinium (: 6.15 eV) is somewhat higher than both the preceding (: 5.64 eV), (: 5.67 eV) and following elements (: 5.86 eV), (: 5.94 eV). This anomaly is due to the fact that gadolinium valence d-subshell borrows 1 electron from the valence f-subshell. Now the valence subshell is the d-subshell, and due to the poor shielding of positive nuclear charge by electrons of the f-subshell, the electron of the valence d-subshell experiences a greater attraction to the nucleus, therefore increasing the energy required to remove the (outermost) valence electron.
Moving into d-block elements: The elements Sc with a 3d1 electronic configuration has a higher IP (: 6.56 eV) than the preceding element (: 6.11 eV), contrary to the decreases on moving into s-block and p-block elements. The 4s and 3d electrons have similar shielding ability: the 3d orbital forms part of the n=3 shell whose average position is closer to the nucleus than the 4s orbital and the n=4 shell, but electrons in s orbitals experience greater penetration into the nucleus than electrons in d orbitals. So the mutual shielding of 3d and 4s electrons is weak, and the effective nuclear charge acting on the ionized electron is relatively large. Yttrium () similarly has a higher IP (6.22 eV) than : 5.69 eV.
Moving into f-block elements; The elements (: 5.18 eV) and (: 5.17 eV) have only very slightly lower IP's than their preceding elements (: 5.21 eV) and (: 5.18 eV), though their atoms are anomalies in that they add a d-electron rather than an f-electron. As can be seen in the above graph for ionization energies, the sharp rise in IE values from (: 3.89 eV) to (: 5.21 eV) is followed by a small increase (with some fluctuations) as the f-block proceeds from to . This is due to the lanthanide contraction (for lanthanides). This decrease in ionic radius is associated with an increase in ionization energy in turn increases, since the two properties correlate to each other. As for d-block elements, the electrons are added in an inner shell, so that no new shells are formed. The shape of the added orbitals prevents them from penetrating to the nucleus so that the electrons occupying them have less shielding capacity.
Ionization energy anomalies in groups
Ionization energy values tend to decrease on going to heavier elements within a group as shielding is provided by more electrons and overall, the valence shells experience a weaker attraction from the nucleus, attributed to the larger covalent radius which increase on going down a group Nonetheless, this is not always the case. As one exception, in Group 10 palladium (: 8.34 eV) has a higher ionization energy than nickel (: 7.64 eV), contrary to the general decrease for the elements from technetium to xenon . Such anomalies are summarized below:
Group 1:
Hydrogen's ionization energy is very high (at 13.59844 eV), compared to the alkali metals. This is due to its single electron (and hence, very small electron cloud), which is close to the nucleus. Likewise, since there are not any other electrons that may cause shielding, that single electron experiences the full net positive charge of the nucleus.
Francium's ionization energy is higher than the precedent alkali metal, cesium. This is due to its (and radium's) small ionic radii owing to relativistic effects. Because of their large mass and size, this means that its electrons are traveling at extremely high speeds, which results in the electrons coming closer to the nucleus than expected, and they are consequently harder to remove (higher IE).
Group 2: Radium's ionization energy is higher than its antecedent alkaline earth metal barium, like francium, which is also due to relativistic effects. The electrons, especially the 1s electrons, experience very high effective nuclear charges. To avoid falling into the nucleus, the 1s electrons must move at very high speeds, which causes the special relativistic corrections to be substantially higher than the approximate classical momenta. By the uncertainty principle, this causes a relativistic contraction of the 1s orbital (and other orbitals with electron density close to the nucleus, especially ns and np orbitals). Hence this causes a cascade of electron changes, which finally results in the outermost electron shells contracting and getting closer to the nucleus.
Group 4:
Hafnium's near similarity in IE with zirconium. The effects of the lanthanide contraction can still be felt after the lanthanides. It can be seen through the former's smaller atomic radius (which contradicts the observed periodic trend ) at 159 pm (empirical value), which differs from the latter's 155 pm. This in turn makes its ionization energies increase by 18 kJ/mol−1.
Titanium's IE is smaller than that of both hafnium and zirconium. Hafnium's ionization energy is similar to zirconium's due to lanthanide contraction. However, why zirconium's ionization energy is higher than the preceding elements' remains unclear; we cannot attribute it to atomic radius as it is higher for zirconium and hafnium by 15 pm. We also cannot invoke the condensed ionization energy, as it is more or less the same ([Ar] 3d2 4s2 for titanium, whereas [Kr] 4d2 5s2 for zirconium). Additionally, there are no half-filled nor fully filled orbitals we might compare. Hence, we can only invoke zirconium's full electron configuration, which is 1s22s22p63s23p63d104s24p64d25s2. The presence of a full 3d-block sublevel is tantamount to a higher shielding efficiency compared to the 4d-block elements (which are only two electrons).
Group 5: akin to Group 4, niobium and tantalum are analogous to each other, due to their electron configuration and to the lanthanide contraction affecting the latter element. Ipso facto, their significant rise in IE compared to the foremost element in the group, vanadium, can be attributed due to their full d-block electrons, in addition to their electron configuration. Another intriguing notion is niobium's half-filled 5s orbital; due to repulsion and exchange energy (in other words the "costs" of putting an electron in a low-energy sublevel to completely fill it instead of putting the electron in a high-energy one) overcoming the energy gap between s- and d-(or f) block electrons, the EC does not follow the Madelung rule.
Group 6: like its forerunners groups 4 and 5, group 6 also record high values when moving downward. Tungsten is once again similar to molybdenum due to their electron configurations. Likewise, it is also attributed to the full 3d-orbital in its electron configuration. Another reason is molybdenum's half filled 4d orbital due to electron pair energies violating the aufbau principle.
Groups 7-12 6th period elements (rhenium, osmium, iridium, platinum, gold and mercury): All of these elements have extremely high ionization energies compared to the elements preceding them in their respective groups. The essence of this is due to the lanthanide contraction's influence on post lanthanides, in addition to the relativistic stabilization of the 6s orbital.
Group 13:
Gallium's IE is higher than aluminum's. This is once again due to d-orbitals, in addition to scandide contraction, providing weak shielding, and hence the effective nuclear charges are augmented.
Thallium's IE, due to poor shielding of 4f electrons in addition to lanthanide contraction, causes its IE to be increased in contrast to its precursor indium.
Group 14: Lead's unusually high ionization energy (: 7.42 eV) is, akin to that of group 13's thallium, a result of the full 5d and 4f subshells. The lanthanide contraction and the inefficient screening of the nucleus by the 4f electrons results in slightly higher ionization energy for lead than for tin (: 7.34 eV).
Bohr model for hydrogen atom
The ionization energy of the hydrogen atom () can be evaluated in the Bohr model, which predicts that the atomic energy level has energy
RH is the Rydberg constant for the hydrogen atom. For hydrogen in the ground state and so that the energy of the atom before ionization is simply
After ionization, the energy is zero for a motionless electron infinitely far from the proton, so that the ionization energy is
. This agrees with the experimental value for the hydrogen atom.
Quantum-mechanical explanation
According to the more complete theory of quantum mechanics, the location of an electron is best described as a probability distribution within an electron cloud, i.e. atomic orbital. The energy can be calculated by integrating over this cloud. The cloud's underlying mathematical representation is the wavefunction, which is built from Slater determinants consisting of molecular spin orbitals. These are related by Pauli's exclusion principle to the antisymmetrized products of the atomic or molecular orbitals.
There are two main ways in which ionization energy is calculated. In general, the computation for the Nth ionization energy requires calculating the energies of and electron systems. Calculating these energies exactly is not possible except for the simplest systems (i.e. hydrogen and hydrogen-like elements), primarily because of difficulties in integrating the electron correlation terms. Therefore, approximation methods are routinely employed, with different methods varying in complexity (computational time) and accuracy compared to empirical data. This has become a well-studied problem and is routinely done in computational chemistry. The second way of calculating ionization energies is mainly used at the lowest level of approximation, where the ionization energy is provided by Koopmans' theorem, which involves the highest occupied molecular orbital or "HOMO" and the lowest unoccupied molecular orbital or "LUMO", and states that the ionization energy of an atom or molecule is equal to the negative value of energy of the orbital from which the electron is ejected. This means that the ionization energy is equal to the negative of HOMO energy, which in a formal equation can be written as:
Molecules: vertical and adiabatic ionization energy
Ionization of molecules often leads to changes in molecular geometry, and two types of (first) ionization energy are defined – adiabatic and vertical.
Adiabatic ionization energy
The adiabatic ionization energy of a molecule is the minimum amount of energy required to remove an electron from a neutral molecule, i.e. the difference between the energy of the vibrational ground state of the neutral species (v" = 0 level) and that of the positive ion (v' = 0). The specific equilibrium geometry of each species does not affect this value.
Vertical ionization energy
Due to the possible changes in molecular geometry that may result from ionization, additional transitions may exist between the vibrational ground state of the neutral species and vibrational excited states of the positive ion. In other words, ionization is accompanied by vibrational excitation. The intensity of such transitions is explained by the Franck–Condon principle, which predicts that the most probable and intense transition corresponds to the vibrationally excited state of the positive ion that has the same geometry as the neutral molecule. This transition is referred to as the "vertical" ionization energy since it is represented by a completely vertical line on a potential energy diagram (see Figure).
For a diatomic molecule, the geometry is defined by the length of a single bond. The removal of an electron from a bonding molecular orbital weakens the bond and increases the bond length. In Figure 1, the lower potential energy curve is for the neutral molecule and the upper surface is for the positive ion. Both curves plot the potential energy as a function of bond length. The horizontal lines correspond to vibrational levels with their associated vibrational wave functions. Since the ion has a weaker bond, it will have a longer bond length. This effect is represented by shifting the minimum of the potential energy curve to the right of the neutral species. The adiabatic ionization is the diagonal transition to the vibrational ground state of the ion. Vertical ionization may involve vibrational excitation of the ionic state and therefore requires greater energy.
In many circumstances, the adiabatic ionization energy is often a more interesting physical quantity since it describes the difference in energy between the two potential energy surfaces. However, due to experimental limitations, the adiabatic ionization energy is often difficult to determine, whereas the vertical detachment energy is easily identifiable and measurable.
Analogs of ionization energy to other systems
While the term ionization energy is largely used only for gas-phase atomic, cationic, or molecular species, there are a number of analogous quantities that consider the amount of energy required to remove an electron from other physical systems.
Electron binding energy
Electron binding energy is a generic term for the minimum energy needed to remove an electron from a particular electron shell for an atom or ion, due to these negatively charged electrons being held in place by the electrostatic pull of the positively charged nucleus. For example, the electron binding energy for removing a 3p3/2 electron from the chloride ion is the minimum amount of energy required to remove an electron from the chlorine atom when it has a charge of −1. In this particular example, the electron binding energy has the same magnitude as the electron affinity for the neutral chlorine atom. In another example, the electron binding energy refers to the minimum amount of energy required to remove an electron from the dicarboxylate dianion −O2C(CH2)8CO.
The graph to the right shows the binding energy for electrons in different shells in neutral atoms. The ionization energy is the lowest binding energy for a particular atom (although these are not all shown in the graph).
Solid surfaces: work function
Work function is the minimum amount of energy required to remove an electron from a solid surface, where the work function for a given surface is defined by the difference
where is the charge of an electron, is the electrostatic potential in the vacuum nearby the surface, and is the Fermi level (electrochemical potential of electrons) inside the material.
Note
| Physical sciences | Periodic table | Chemistry |
59615 | https://en.wikipedia.org/wiki/Electric%20potential | Electric potential | Electric potential (also called the electric field potential, potential drop, the electrostatic potential) is defined as the amount of work/energy needed per unit of electric charge to move the charge from a reference point to a specific point in an electric field. More precisely, the electric potential is the energy per unit charge for a test charge that is so small that the disturbance of the field under consideration is negligible. The motion across the field is supposed to proceed with negligible acceleration, so as to avoid the test charge acquiring kinetic energy or producing radiation. By definition, the electric potential at the reference point is zero units. Typically, the reference point is earth or a point at infinity, although any point can be used.
In classical electrostatics, the electrostatic field is a vector quantity expressed as the gradient of the electrostatic potential, which is a scalar quantity denoted by or occasionally , equal to the electric potential energy of any charged particle at any location (measured in joules) divided by the charge of that particle (measured in coulombs). By dividing out the charge on the particle a quotient is obtained that is a property of the electric field itself. In short, an electric potential is the electric potential energy per unit charge.
This value can be calculated in either a static (time-invariant) or a dynamic (time-varying) electric field at a specific time with the unit joules per coulomb (J⋅C−1) or volt (V). The electric potential at infinity is assumed to be zero.
In electrodynamics, when time-varying fields are present, the electric field cannot be expressed only as a scalar potential. Instead, the electric field can be expressed as both the scalar electric potential and the magnetic vector potential. The electric potential and the magnetic vector potential together form a four-vector, so that the two kinds of potential are mixed under Lorentz transformations.
Practically, the electric potential is a continuous function in all space, because a spatial derivative of a discontinuous electric potential yields an electric field of impossibly infinite magnitude. Notably, the electric potential due to an idealized point charge (proportional to , with the distance from the point charge) is continuous in all space except at the location of the point charge. Though electric field is not continuous across an idealized surface charge, it is not infinite at any point. Therefore, the electric potential is continuous across an idealized surface charge. Additionally, an idealized line of charge has electric potential (proportional to , with the radial distance from the line of charge) is continuous everywhere except on the line of charge.
Introduction
Classical mechanics explores concepts such as force, energy, and potential. Force and potential energy are directly related. A net force acting on any object will cause it to accelerate. As an object moves in the direction of a force acting on it, its potential energy decreases. For example, the gravitational potential energy of a cannonball at the top of a hill is greater than at the base of the hill. As it rolls downhill, its potential energy decreases and is being translated to motion – kinetic energy.
It is possible to define the potential of certain force fields so that the potential energy of an object in that field depends only on the position of the object with respect to the field. Two such force fields are a gravitational field and an electric field (in the absence of time-varying magnetic fields). Such fields affect objects because of the intrinsic properties (e.g., mass or charge) and positions of the objects.
An object may possess a property known as electric charge. Since an electric field exerts force on a charged object, if the object has a positive charge, the force will be in the direction of the electric field vector at the location of the charge; if the charge is negative, the force will be in the opposite direction.
The magnitude of force is given by the quantity of the charge multiplied by the magnitude of the electric field vector,
Electrostatics
An electric potential at a point in a static electric field is given by the line integral
where is an arbitrary path from some fixed reference point to ; it is uniquely determined up to a constant that is added or subtracted from the integral. In electrostatics, the Maxwell-Faraday equation reveals that the curl is zero, making the electric field conservative. Thus, the line integral above does not depend on the specific path chosen but only on its endpoints, making well-defined everywhere. The gradient theorem then allows us to write:
This states that the electric field points "downhill" towards lower voltages. By Gauss's law, the potential can also be found to satisfy Poisson's equation:
where is the total charge density and denotes the divergence.
The concept of electric potential is closely linked with potential energy. A test charge, , has an electric potential energy, , given by
The potential energy and hence, also the electric potential, is only defined up to an additive constant: one must arbitrarily choose a position where the potential energy and the electric potential are zero.
These equations cannot be used if i.e., in the case of a non-conservative electric field (caused by a changing magnetic field; see Maxwell's equations). The generalization of electric potential to this case is described in the section .
Electric potential due to a point charge
The electric potential arising from a point charge, , at a distance, , from the location of is observed to be
where is the permittivity of vacuum, is known as the Coulomb potential. Note that, in contrast to the magnitude of an electric field due to a point charge, the electric potential scales respective to the reciprocal of the radius, rather than the radius squared.
The electric potential at any location, , in a system of point charges is equal to the sum of the individual electric potentials due to every point charge in the system. This fact simplifies calculations significantly, because addition of potential (scalar) fields is much easier than addition of the electric (vector) fields. Specifically, the potential of a set of discrete point charges at points becomes
where
is a point at which the potential is evaluated;
is a point at which there is a nonzero charge; and
is the charge at the point .
And the potential of a continuous charge distribution becomes
where
is a point at which the potential is evaluated;
is a region containing all the points at which the charge density is nonzero;
is a point inside ; and
is the charge density at the point .
The equations given above for the electric potential (and all the equations used here) are in the forms required by SI units. In some other (less common) systems of units, such as CGS-Gaussian, many of these equations would be altered.
Generalization to electrodynamics
When time-varying magnetic fields are present (which is true whenever there are time-varying electric fields and vice versa), it is not possible to describe the electric field simply as a scalar potential because the electric field is no longer conservative: is path-dependent because (due to the Maxwell-Faraday equation).
Instead, one can still define a scalar potential by also including the magnetic vector potential . In particular, is defined to satisfy:
where is the magnetic field. By the fundamental theorem of vector calculus, such an can always be found, since the divergence of the magnetic field is always zero due to the absence of magnetic monopoles. Now, the quantity
is a conservative field, since the curl of is canceled by the curl of according to the Maxwell–Faraday equation. One can therefore write
where is the scalar potential defined by the conservative field .
The electrostatic potential is simply the special case of this definition where is time-invariant. On the other hand, for time-varying fields,
unlike electrostatics.
Gauge freedom
The electrostatic potential could have any constant added to it without affecting the electric field. In electrodynamics, the electric potential has infinitely many degrees of freedom. For any (possibly time-varying or space-varying) scalar field, , we can perform the following gauge transformation to find a new set of potentials that produce exactly the same electric and magnetic fields:
Given different choices of gauge, the electric potential could have quite different properties. In the Coulomb gauge, the electric potential is given by Poisson's equation
just like in electrostatics. However, in the Lorenz gauge, the electric potential is a retarded potential that propagates at the speed of light and is the solution to an inhomogeneous wave equation:
Units
The SI derived unit of electric potential is the volt (in honor of Alessandro Volta), denoted as V, which is why the electric potential difference between two points in space is known as a voltage. Older units are rarely used today. Variants of the centimetre–gram–second system of units included a number of different units for electric potential, including the abvolt and the statvolt.
Galvani potential versus electrochemical potential
Inside metals (and other solids and liquids), the energy of an electron is affected not only by the electric potential, but also by the specific atomic environment that it is in. When a voltmeter is connected between two different types of metal, it measures the potential difference corrected for the different atomic environments. The quantity measured by a voltmeter is called electrochemical potential or fermi level, while the pure unadjusted electric potential, , is sometimes called the Galvani potential, . The terms "voltage" and "electric potential" are a bit ambiguous but one may refer to of these in different contexts.
Common formulas
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59637 | https://en.wikipedia.org/wiki/Crocus | Crocus | Crocus (; plural: crocuses or croci) is a genus of seasonal flowering plants in the family Iridaceae (iris family) comprising about 100 species of perennials growing from corms. They are low growing plants, whose flower stems remain underground, that bear relatively large white, yellow, orange or purple flowers and then become dormant after flowering. Many are cultivated for their flowers, appearing in autumn, winter, or spring. The flowers close at night and in overcast weather conditions. The crocus has been known throughout recorded history, mainly as the source of saffron. Saffron is obtained from the dried stigma of Crocus sativus, an autumn-blooming species. It is valued as a spice and dyestuff, and is one of the most expensive spices in the world. Iran is the center of saffron production. Crocuses are native to woodland, scrub, and meadows from sea level to alpine tundra from the Mediterranean, through North Africa, central and southern Europe, the islands of the Aegean, the Middle East and across Central Asia to Xinjiang in western China. Crocuses may be propagated from seed or from daughter cormels formed on the corm, that eventually produce mature plants. They arrived in Europe from Turkey in the 16th century and became valued as an ornamental flowering plant.
Description
General
Crocus display the general characteristics of family Iridaceae, which include basal cauline (arising from the aerial stem) leaves that sheath the stem base, hermaphrodite flowers that are relatively large and showy, the perianth petaloid with two whorls of three tepals each and septal nectaries. The flowers have three stamens and a gynoecium of three united carpels and an inferior ovary, three locules and axile placentation with fruit that is a loculicidal capsule.
Crocus is an acaulescent (lacking a visible lower stem above ground) diminutive seasonal cormous (growing from corms) herbaceous perennial geophytic genus. The corms are symmetrical and globose or oblate (round in shape with flatted tops and bottoms), and are covered with tunic leaves that are fibrous, membranous or coriaceous (leathery). The corms produce fibrous roots, and contractile roots which adjust the corms depth in the soil, which may be pulled as deep as into the soil.The roots appear randomly from the lower part of the corm, but in a few species, from a basal ridge.
Leaves
Plants produce several basal linear bifacial green leaves that arise from the corms. These are adaxially (upper surface facing axis) flat or channelled with pale median stripes, while the opposite (abaxial) surface is strongly keeled, with two grooves on either side. The leaves have a distinctive shape in cross-section, being boat-shaped with two lateral arms with margins recurved inwardly towards the central ridged keel, forming the sides of the "boat". The keel may be square or rectangular, but is lacking in C. carpetanus. The pale central stripe is caused by parenchymatous cells which lack chloroplasts and may contain air spaces. The leaves are from wide and long. The leaf-like bracts are membranous, while the smaller bracteoles are either membranous or absent. The leaf bases are surrounded by up to 5 membranous sheaths called cataphylls, a specialised leaf. The bases of the cataphylls form the corm tunic, and their number varies from 3 to 6, and enclose the true leaves (euphylls), bracts, bracteoles and flowering stalk.
Flowers
The number of peduncles (flower stems) vary from one to several and remain underground, emerging only at the fruiting stage, bearing flowers that are solitary or several, so that a true scape is absent. The flowers are pedicellate (attached to the peduncle by a short subterranean pedicel stalk). The pedicel is sometimes subtended (below pedicel) by a membranous, sheathing prophyll (leaf-like structure).
The showy, salver to cup-shaped, single or clustered actinomorphic flowers taper off into a narrow tube; the flowers emerge from the ground, and can be white, yellow, lilac to dark purple, or variegated in cultivars. The flower tube is long, cylindrical and slender, expanding apically. The floral tube is long and narrow with 6 lobes in 2 whorls. The perianth is 3+3 (3 sepals+3petals) and gamophyllous (with fused segments). The tepal whorls are similar, equal or subequal with a smaller inner whorl, and cupped to outspread. The bracts are membranous, but the inner ones are sometimes lacking.
The 3 stamens are erect and linear and inserted in the throat of the perianth tube, with anthers shorter than the filaments. Pollen grains are inaperturate (apertures absent) but sometimes spiraperturate (spiral shaped). Each flower has a single style which is exserted (projecting beyond the corolla tube) and slender distally with three to many branches. The branches are highly variable, being short or long, and simple, bifurcate (dividing in two) or multifid and sometimes distally flattened. The inferior ovary has 3 carpels with axile placentation. It remains underground, and as the seeds ripen, the pedicel (stem of the flower) grows longer - so the fruit is above the soil surface.
Fruit and seed
The fruit is a small membranous capsule, ellipsoid or oblong-ellipsoid in shape and the many seeds are globose to ellipsoid. The seed surface is highly variable, including papillate (covered in small protuberances), digitiform (finger-like) and other epidermal cell types. In some species the seeds are arillate, with fleshy appendages. Crocus seeds have both inner and outer integuments and in some species the outer epidermis may display long papillae. Embryo-sac development is Polygonum type. Dehiscence (splitting of the capsule to release the seed) is of the loculicidal type in which it splits through the wall of the locules leaving the septa that separate them intact.
Karyology
Crocus has extensive aneuploidy (abnormal number of chromosomes), with some uncertainty as to the base number of chromosomes. The chromosome numbers shows extreme variability, ranging from 2n=6 to 2n=70 even within a single species.
Phytochemistry
The Iridaceae contain a wide range of phenolic compounds. However, 6-Hydroxyflavones are found only in Crocus, which is also characterised by the presence of crocins, water-soluble yellow carotenoids, in the floral tissues. Crocin is a diester of crocetin, responsible for the colour of the styles and stigma of C. sativus, and hence saffron. A few species contain mangiferin, a .
While the flowers may vary dramatically between species, there is little variation in the leaves, but sufficient variability in corm tunics that they may be used as an aid in differentiating taxa.
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Taxonomy
History
The crocus was well known to the ancients, being described at least as early as Theophrastus (c. 371c. 287 BC), and was introduced into Britain by the Romans, where the saffron crocus was used as a dyestuff. It was reintroduced into Western Europe by the Crusaders. The crocus is mentioned in mediaeval and later herbals, one of the earliest being the 14th century Tractatus de Herbis. William Turner (1548) states that the crocus is referred to as saffron in English, implying that only C. sativus was known at that time. However, by 1597 John Gerard writes of "sundry sorts" and uses the term saffron and crocus as interchangeable. He included both spring and autumn flowering crocus, but distinguished Wild Saffron (Crocus) from Meadow Saffron (Colchicum). He described eleven forms. Some of his specimens were obtained from Clusius. In the following century, John Parkinson in a more detailed account was more careful to include separate chapters for Colchicum, with the common name of meadow saffron, from Crocus or saffron. Parkinson (1656) states that there are "divers sorts of saffrons" describing 27 spring flowering plants and 4 autumn flowering ones, pointing out that only one of those was the true saffron crocus, which he called Crocus verus sativus autumnalis. Similar accounts are found in continental European herbals, including those of l'Obel in Flanders (1576) and Besler's Hortus Eystettensis in Bavaria (1613).
The genus Crocus was first formally described by Linnaeus in 1753, with three taxa, and two species, C. sativus (type species), var. officinalis (now treated as a synonym of C. sativus) and var. vernus (now C. vernus) and C. bulbocodium (now Romulea bulbocodium). Thus Linnaeus recognised two taxa that are accepted as separate species in modern classifications, one vernal and one autumnal crocus, but incorrectly assumed they were only varieties of a single species, while his second species was actually from a closely related genus that was only recognised later (1772). However, a subsequent re-examination of Linnaeus's specimens suggested the presence of several different species that he did not recognise as being separate. Linnaeus' system, based on sexual characteristics, Crocus was classified as Triandra Monogynia (Three stamens, Single pistil). Linnaeus's system was supplanted by the "natural" system which used a hierarchy of taxonomic ranks based on weighting of the importance of structural characteristics of the plant. Jussieu (1789) placed the genus Crocus in his Ordo (family) Irides or Les iris, as a member of the class Stamina epigyna (stamens inserted above the ovary) as part of the monocotyledons, the first level of the division of the flowering plants.
One of the first monographs of the genus appeared in 1809, by Haworth, followed in 1829 by that of Sabine, and Herbert in 1847. In 1853, Lindley continued the placement of Crocus as one of 53 genera in Iridaceae, which he included in a higher order of monocotyledons, the Narcissales. Baker published a monograph on the genus in 1874, adopting a very different schema to that of Herbert. In 1883, Bentham and Hooker described the Irideae (Iridaceae) as having more than 700 species, and divided it into 3 tribes and further into subtribes. Tribe Sysyrinchieae as having 2 subtribes, including Ixieae. The latter was circumscribed with four genera, Crocus, Syringodea, Galaxia (Moraea) and Romulea. This circumscription has remained stable since, with the exception of Moraea which properly belongs in a separate tribe. The most influential monograph of the nineteenth century was that of Maw (1886), which forms the basis of modern understanding of the genus. Maw built on the work of Herbert, rejecting Baker's classification. The availability of molecular phylogenetic methods in the late twentieth century has shown that the Iridaceae properly belong within the order Asparagales.
Botanical illustration
The scientific study of the genus in the late eighteenth century was accompanied by detailed descriptions with Botanical illustrations, such as those of William Curtis (1787) and Sims (1803), that appeared in Curtis' Botanical Magazine, with illustrations by Sydenham Edwards. Other illustrations are found in monographs such as those of Haworth (1809) and Sabine (1830), illustrated by Charles John Robertson. The largest collection is found in the most comprehensive monograph, that of Maw (1886). Other sources include the portfolios of plates, such as the survey of the plants of France by Masclef (1891). At that time only C. sativus and C. vernus were included in the Flora of France.
Phylogeny
The genus Crocus belongs to the monocot family Iridaceae (iris family), specifically the large subfamily Crocoideae. Within that subfamily, crocus is placed on the tribe, Ixieae (synonym Croceae), one of five. The Ixieae are then subdivided into subtribes, with the genera Crocus, Romulea and Syringodea forming subtribe Romuleinae. The Romuleinae have been characterised within the Ixieae by progressively reduced aerial stems. solitary flowers on the stem branches and woody tunics on the corms. They also often have divided style branches. However, Crocus corm tunics are fibrous and membranous rather than woody as in Syringodea. Also, Crocus has a ridged and often keeled abaxial leaf surface, while that of Syringodea is rounded, and the midline adaxial translucency of Crocus is lacking in Syringodea. Romulea is principally distinguished from the other two genera by generally having aerial stems or at least an ovary at ground level, compared with the other acaulescent genera, other differences include unifacial rather than bifacial leaves and the pollen structure.
Within the Romuleinae, Crocus is a sister group to Syringodea, the two genera forming a sister group to Romulea.
Subdivision
The genus Crocus consists of about 200 accepted species, which continue to increase, and has undergone a large number of taxonomic classifications. The genus has often been divided into sections, beginning with that of Haworth (1809) who described two sections based on the presence or absence of hairs in the throat of the flower, while Sabine was the first to realise the importance of the presence or absence of a basal spathe (prophyll) in dividing the genus into two sections, a practice followed by Herbert. However, Sabine's practice of using trinomials for varieties such as C. sulphureus concolor is no longer accepted, although Herbert somewhat similarly used varieties and subvarieties, e.g. C. vernus var.1 Communis subvar. 1. Obovatus. Herbert also used geographical distribution as a basis of classification. By the late 19th century Maw (1886), following Herbert, subdivided the genus into two divisions, the Involucrati and the Nudiflori, and then further divided it into six sections and lastly by flowering times (spring or autumn). Although rejecting the concept of subvarieties, he placed even more emphasis on geography.
The most widely accepted system, that proposed by Brian Mathew in 1982 was based on Maw's system, but with less emphasis on flowering times. This mainly depended on three character states:
the presence or absence of a prophyll (a basal spathe);
the aspect of the style;
the corm tunic.
and included 81 species, however, one of these, Crocus medius was later recognized as a synonym of Crocus nudiflorus.
The genus, as described by Mathew, consisted of two subgenera, Crocirus (monotypic for Crocus banaticus) and Crocus including the remainder of the species, based on whether the anthers were introrse or extrorse (dehiscence directed towards or away from centre of flower) respectively. Subgenus Crocus was then divided into two sections, Crocus and Nudiscapus, based on the presence or absence of the prophyll. Each section was then further divided into six series of Crocus and nine of Nudiscapus. These series were defined by the division of the style, the corm tunic, flowering time, leaf structure, presence of a bracteole and anther colour. Mathew also introduced the concept of subspecies, including 50 in all, by giving similar but different forms subspecies status if geographically separated, resulting in about 140 distinct taxa. The seven species and ten subspecies discovered since then have been integrated into revisions of this classification, though new species continue to be described, leading to estimates of at least 200 species.
Speciation
Crocus populations have extremely high infra-specific variability with a very diverse spectrum of morphological and phenotypical varieties, while many individual specimens from different species may closely resemble each other. Based on such morphological differences between isolated populations many new species have been named, but without a definition of new species based on molecular and/or karyological information, species can not be confirmed, creating difficulties in determining speciation and hence the exact number of species. The situation is even more complex once hybridisation (combination of taxa) and introgression (transfer of genetic material) are considered.
Molecular phylogeny
The availability of molecular phylogeny methods revealed problems with the traditional systems based on morphology alone. The first analysis of the complete genus was carried out by Mathew and colleagues in 2008 using nucleotide sequences from plastid regions. In particular, the DNA data suggest there are no grounds for isolating C. banaticus in its own subgenus Crociris, though it is a unique species in the genus. Because it has a prophyll at the base of the pedicel, it therefore would fall within section Crocus, although its exact relationship to the rest of the subgenus remains unclear.
Of the 15 series in the Mathew scheme, only seven were monophyletic, and in particular the largest series, Biflori and Reticulati, which include a third of all species, were non-monophyletic. Another anomalous species, C. baytopiorum, should now be placed in a series of its own, series Baytopi. C. gargaricus subsp. herbertii has been raised to species status, as C. herbertii. The autumn-flowering C. longiflorus, the type species of series Longiflori (long regarded by Mathew as "a disparate assemblage"), appeared to lie within series Verni. In addition, the position of C. malyi was currently unclear.
DNA analysis and morphological studies suggest further that series Reticulati, Biflori and Speciosi are "probably inseparable", C. adanensis and C. caspius should probably be removed from Biflori, C. adanensis falls in a clade with C. paschei as a sister group to the species of series Flavi and C. caspius appears to be sister to the species of series Orientales.
The study showed "no support for a system of sections as currently defined", although, despite the many inconsistencies between Mathew's 1982 classification and the current hypothesis, "the main assignment of species to the sections and series of that system is actually supported". The authors state, "further studies are required before any firm decisions about a hierarchical system of classification can be considered" and conclude "future re-classification is likely to involve all infrageneric levels, subgenera, sections and series". A further study, using the internal transcribed spacer region (ITS) of the nuclear ribosomal DNA (rDNA), together with a chloroplast marker, broadly confirmed these findings.
Crocus forms a monophyletic clade, with a basal polytomy of four subclades. The first clade (A) corresponding to section Crocus, but including C. sieberi and several closely related species (originally included in section Nudiscapus series Reticulati). The remaining three clades (B-D) include all the remaining species of section Nudiscapus. Of these, B and C are small, corresponding to series Orientales and Carpetani respectively, with all remaining series in the large D clade. The exception is C. caspius, originally in series Biflori, which segregates in clade B. Thus, although division of the genus into two sections is well supported, no single morphological character defines these two groups. The C. sieberi group are assumed to have lost their prophyll secondarily. Of the series, eight could be shown to be monophyletic; Crocus, Kotschyani and Scardici (section Crocus) and Aleppici, Carpetani, Laevigati, Orientalis and Speciosi (section Nudiscapus). Flowering season did not correspond to molecular groupings and nor did any of the previously used morphological characteristics, indicating a high degree of homoplasy, in which traits are gained or lost independently in different lineages. The remainder of the series could not be supported as natural groupings. Mathew's concept of subspecies status within C. biflorus could not be supported, each being considered a separate species, resulting in the genus having at least 150 species.
A more detailed molecular and morphological study of series Verni (section Crocus) allowed it to be better characterised and circumscribed, as well as the closely related series Longiflori. Series Verni sensu Mathew was found to consist of two groups, the first being C. vernus sensu Mathew and the other consisting of C. etruscus, ilvensis, kosaninii and longiflorus. The taxonomic status of C. vernus had been uncertain for some time, given the observation that the name was more properly applied to C. albiflorus, requiring a new designation of C. neapolitanus for those previously known as C. vernus. Subsequently C. vernus was split into 5 separate species. The incorporation of C. longiflorus into series Verni resulted in making series Longiflori no longer a legitimate taxonomic unit.
In section Nudiscapus, series Reticulati was polyphyletic with species intermingled with series Biflori and Speciosi, requiring a recircumscription, confining Reticulati to 8 species, to obtain monophyly. Among the thereby displaced species, are a number of very closely related taxa, referred to as the Crocus sieberi aggregate, which has been proposed as a new series Sieberi. Other new series, such as Isauri and Lyciotauri, continue to be created out of the Biflori series.
Mathew's circumscription of Crocus introduced the rank of subspecies, of which the largest number (14) were those of Crocus biflorus Miller, the type species of series Biflori, a number which continued to grow. Molecular methods identified these as a polyphyletic assemblage rather than closely related subordinate infraspecific taxa. This necessitated a complete taxonomic revision of series Biflori, elevating each subspecies to species status. A similar issue occurs with C. reticulatus sensu Mathew, who created two subspecies, resulting in 9 newly defined species.
Sections and species
The classification of Brian Mathew (1982), as amended in 2009 divides the genus into two sections, further divided by series. The number of series, continues to evolve.
Section Crocus B.Mathew
Species with a basal prophyll. Type species C. sativus L.
6 series
Section Nudiscapus B.Mathew
Species without a basal prophyll. Type species C. reticulatus Stev. ex Adams
9 series
Similarly named species
Some crocus species, known as "autumn crocus", flower in late summer and autumn, during (autumnal) rains, after summer's heat and drought. The name autumn crocus is also often used as a common name for Colchicum, which is not a true crocus but in its own family (Colchicaceae) in the lily order Liliales. The plants are toxic, but have medicinal uses. Colchicum are also known as meadow saffron, though true saffron is not toxic. Crocus species have three stamens while Colchicum species have six; crocus have one style, while Colchicum have three.
Some Pulsatilla species are also called "prairie crocus" (previously Anemone patens) or "wild crocus", but they belong to the buttercup family (Ranunculaceae). Pulsatilla species, which are commonly called pasqueflowers, unlike crocuses have rhizomes, the foliage is covered with long soft hairs, and the flowers are produced on above-ground stems.
Etymology
"Crocus", the name of the genus, is Late Middle English (late 14th century) and also denotes saffron. It is derived via Latin crocus from the Greek κρόκος (krokos), which is itself probably a loan word from a Semitic language, related to Hebrew כרכום karkōm, Aramaic ܟܟܘܪܟܟܡܡܐ kurkama, and Arabic kurkum, meaning saffron (Crocus sativus), "saffron yellow" or turmeric (see Curcuma), another yellow dye. The word ultimately traces back to the Sanskrit kunkumam () for "saffron". The English name is a learned 16th-century adoption from the Latin safranum, but Old English already had croh for saffron, introduced by the Romans.
Distribution and habitat
Crocuses are distributed from the Mediterranean, from the Iberian peninsula and North Africa, through central and southern Europe, the islands of the Aegean, the Middle East and across central and southwest Asia to Xinjiang in western China, but most species are restricted to Turkey and Asia Minor and the Balkans, with the Balkan Peninsula having the largest number of species (at least 31), forming the centre of diversity, however they are widely introduced. The distribution of species is described over five contiguous areas from west to east (see map).
Habitats range from sea level to as high as subalpine altitudes, and in a wide range of habitats from woodlands to meadows and deserts, often on stony mountain slopes with good drainage. The majority of species are native to areas with cold winters and hot summers with little rain, and active growth is typically from fall to mid-spring. The natural habitats of crocus species are threatened by human activities, including urbanization, industrialization, and other land disturbances and recreational uses. They are negatively impacted by uncontrolled gathering and heavy grazing by livestock.
Ecology
The life cycle of Crocus species begins with the seed, germinating to a seedling, and a mature plant in 3–5 years, however seeds may remain dormant in the soil for several years. The germination stages was first described and illustrated by Maw in his 1886 monograph. In its first year, the crocus produces only a single leaf and creates a corm covered by a thin tunic, about 5–8 mm in size, dependent on the species. In the northern hemisphere, the autumnal crocuses flower between September and November. The vernal (spring) crocuses flowering time depends both on climate and habitat, but is usually mid-winter to spring. Leaves may be synanthous (produced during flowering) or hysteranthous (when the flowers wither away). In the summer, with hot and dry conditions the plant becomes dormant, with all the above ground parts dying back. Colder temperatures in winter then activate the corms. Propagation occurs sexually by seed and asexually by small corms, called cormels or cormlets, produced in the axils of the corms (between tunic scales and body of corm). As the fruit capsule ripens, it emerges from the soil at the base of the flowering stem before dehiscing (splitting open) and releasing the seeds. Seed dispersal may be enhanced by ants, at least in species with arillate seeds.
At night and in overcast weather, the perianth closes. The ovary produces nectar which attracts bees (particularly female bumblebees) and Lepidoptera.
Pests and diseases
Cultivated plants may have their corms consumed by mice and other rodents, including voles, squirrels, and chipmunks. They are also attacked by mildew, gray mold, botrytis, and fusarium rot. Root rot may also occur, caused by Stromatinia gladioli and Pythium species - the nematode Pratylenchus penetrans may also cause root rot. Viruses that are known to infect Crocus spp include: Potyviruses, especially bean yellow mosaic virus and also tobacco rattle virus, tobaccos necrosis virus, and cucumber mosaic virus. The foliage may experience rot, rust, and scab diseases and be fed upon by aphids, mites, snails, and slugs. The foliage is eaten by hares, rabbits, and deer; the flowers are sometimes removed by birds, including crows, jackdaws, and magpies.
Cultivation
Saffron
The economic importance of the genus is largely dependent on the single species, Crocus sativus, now known only in cultivation. C. sativus is grown for the production of saffron, an orange-red derivative of its dried stigma, and among the most expensive spices in the world. The estimated worldwide production of C. sativus plants is 205 tons. About 180,000 stigmas from 60,000 flowers are required to produce saffron, which sells for about US$10,000 (2018). Modern saffron production is widely cultivated in Kashmir, Iran, Turkey, Afghanistan and the Mediterranean from Spain to Asia Minor. An important center is the eponymous town of Krokos, in the Kozani region of Greece. The saffron product, Krokos Kozanis is a PDO (Protected Designation of Origin). Production is largely indigenous and Iran accounts for 65% of global production, covering 72,162 ha.
Saffron is thought to have been used in embalming in Ancient Egypt. It is mentioned in the Old Testament, in the Song of Songs as a precious spice and has featured as a dye and fragrance throughout written history, with mention in the Iliad.
Cultivation and harvesting of C. sativus for saffron was first documented in the Mediterranean, notably on the island of Crete. Frescos showing them are found there at the Bronze Age Minoan site of Knossos, as well as from the comparably aged Akrotiri site on the Aegean island of Santorini, and formed an important part of the Minoan economy and culture and had both a sacred role and use as a psychoactive drug and food additive. Women still gather crocuses in the Akrotiri region.
Horticulture and floriculture
Crocuses were described in Turkish gardens in the early sixteenth century, gathered from the far reaches of the Ottoman Empire, where they were seen by visiting European botanists and explorers, among the first of whom was Pierre Belon who arrived in Constantinople in 1547. The first crocus seen in the Netherlands, where crocus species were not native, were from corms brought to Vienna in 1562 from Constantinople by the Holy Roman Emperor's ambassador to the Sublime Porte, Ogier Ghiselin de Busbecq. A few corms were forwarded to Carolus Clusius at the botanical garden in Leiden. These were almost certainly cultivated varieties rather than wild species. European visitors to Turkey continued to bring back specimens for gardens in their own country. Prominent among the latter were the gardens at Middelburg in the Netherlands. Jehan Somer, a Middelburg merchant, brought back crocuses among his other specimens in 1592, where they attracted the attention not only of Clusius but of the early Dutch flower painters, notably Ambrosius Bosschaert. By 1620, new garden varieties had been developed, and featured in contemporary illustrations, such as that of Crispijn van de Passe in his Hortus floridus of 1614. There are accounts of crocus gardens in the seventeenth century, such as the Saffron Garth of Walter Stonehouse at Darfield, Yorkshire.
Crocuses are among the most important ornamental geophytes in the global flower industry, ranking sixth in terms of Dutch bulb production (2003–2008) with 463–668 ha under cultivation. The crocus is one of the most popular flowers found in the garden in the late winter and early spring. About 30 of the species are cultivated, among the most popular being C. chrysanthus, C. flavus, C. sieberi, C. tommasinianus and C. vernus, together with hundreds of cultivars derived from them. Both fall and spring blooming crocuses are cultivated for their flowers. Among the first flowers to bloom in spring, their flowering time can vary from fall to the late winter blooming C. tommasinianus; the earliest fall blooming species, C. scharojanii, may flower during the last weeks of July.
The varieties cultivated for decoration in gardens and pots mainly represent six species: C. vernus, C. chrysanthus, C. flavus, C. sieberi, C. speciosus and C. tommasinianus. During the horticulture production year 2009/2010, more than 70 cultivars were grown in Holland, covering an area of 366 hectares; the most common ones were 'Flower Record' and 'King of the Stripes' which accounted for 42 hectares, other species grown included C. chrysanthus, C. tommasinianus, and C. flavus - all are spring blooming plants. But the most commonly grown plants are the Dutch hybrids with large flowers in a rich palette of colors.
Both sexual and asexual means are used to increase the number of plants; seeds and multiplication of corms are the most common means of production, but tissue culture can be used, most commonly for saffron crocus. New corms are formed on top of the older corm which withers away, and cormels are produced from axillary buds. The production of new plants begins with harvested corms in late June to early July, after being graded by corm size the corms are stored around 22 Celsius until early October when they are moved to 17 Celsius until planted later in October and November; flowering occurs in March and the flowers are not removed. Crocuses are also forced to produce flowering plants out of season and the most common species used are C. vernus and C. flavus, and most of the corms used for forcing come from the Netherlands.
Spring flowering types are planted in fall, while fall-blooming types in late summer; typically, the corms are placed 3 to 4 inches deep in well-draining soil in areas with full sun exposure. They do not thrive in heavy clay soils or those that are damp, especially during their summer dormancy period. Commercial crops are produced on raised beds and slopes, to ensure adequate drainage, while horticulturalists often plant on sand beds for the same purpose. Spring flowering types also do well in areas with deciduous trees, where they flower and produce leaves before the trees completely leaf-out. Crocuses are grown in USDA winter zones 3–8. Not all species are hardy in the upper zones; C. sativus is winter hardy in USDA zones 6 through 8, and C. pulchellus is hardy in zones 5 through 8.
Some are suitable for naturalizing in grass, but mowing off the foliage before it turns yellow produces short lived plants. Some crocuses, especially C. tommasinianus and its selected forms and hybrids (such as 'Whitewell Purple' and 'Ruby Giant'), seed prolifically and are ideal for naturalizing. They can, however, become weeds in rock gardens, where they will often appear in the middle of choice, mat-forming alpine plants, and can be difficult to remove. Crocus flowers and leaves are protected from frost by a waxy cuticle; in areas where snow and frost occasionally occur in the early spring, it is not uncommon to see early flowering crocuses blooming through a light late snowfall.
Autumn crocus
Autumn-flowering species of crocus that are cultivated include:
C. laevigatus has a long flowering period which starts in late autumn or early winter and may continue into February.
Colchicum autumnale is commonly known as "autumn crocus", but is a member of the plant family Colchicaceae, and not a true crocus (of the family Iridaceae).
Uses
The corms of crocuses have been used as foodstuffs in Syria. The carotenoids found in the styles of Crocus species, particularly C. sativus have been shown to inhibit cancer cell proliferation, and have led to interest in potential pharmaceutical applications.
Culture
The crocus or krokos has been known since ancient times, and used in decorative arts, such as the Minoan wall paintings in Santorini from ca. 1,600 BC. Representations of the saffron crocus appear frequently in Minoan art and pervade Aegean art from the Early Bronze Age to the Mycenaean period. Theophrastos (4th century BC) described the saffron crocus as being valued as a spice and dye, while Homer compares a sunrise to the flower colour. Saffron coloured robes were much admired by women in antiquity and gave the garment Crocota its name. The oil was also valued as a cosmetic. According to Greek legend Crocus or Krokus (), was a mortal youth the gods turned into a plant bearing his name, the crocus, after his death caused by his great desire and unfulfilled love for the shepherdess Smilax. Other versions state that as he died three tears fell into the flower becoming its three stigmata.
Crocuses occur in many flower paintings, one of the earliest being that of Ambrosius Bosschaert's Composed Bouquet of Spring Flowers (1620). In this painting the cream-colored crocus feathered with bronze at the base of the bouquet reflected varieties on the market at that time. Bosschaert, working from a preparatory drawing to paint his composed piece spanning the whole of spring, exaggerated the crocus so that it passes for a tulip, but its narrow, grass-like leaves give it away.
The crocus is used in many contexts to symbolically denote spring and new beginnings. For instance, it was used as the emblem of the 2019 FIFA U-20 World Cup in Poland to symbolise the emergence of new talent.
| Biology and health sciences | Asparagales | Plants |
59644 | https://en.wikipedia.org/wiki/Digital%20signature | Digital signature | A digital signature is a mathematical scheme for verifying the authenticity of digital messages or documents. A valid digital signature on a message gives a recipient confidence that the message came from a sender known to the recipient.
Digital signatures are a standard element of most cryptographic protocol suites, and are commonly used for software distribution, financial transactions, contract management software, and in other cases where it is important to detect forgery or tampering.
Digital signatures are often used to implement electronic signatures, which include any electronic data that carries the intent of a signature, but not all electronic signatures use digital signatures. Electronic signatures have legal significance in some countries, including Brazil, Canada, South Africa, Russia, the United States, Algeria, Turkey, India, Indonesia, Mexico, Saudi Arabia, Uruguay, Switzerland, Chile and the countries of the European Union.
Digital signatures employ asymmetric cryptography. In many instances, they provide a layer of validation and security to messages sent through a non-secure channel: Properly implemented, a digital signature gives the receiver reason to believe the message was sent by the claimed sender. Digital signatures are equivalent to traditional handwritten signatures in many respects, but properly implemented digital signatures are more difficult to forge than the handwritten type. Digital signature schemes, in the sense used here, are cryptographically based, and must be implemented properly to be effective. They can also provide non-repudiation, meaning that the signer cannot successfully claim they did not sign a message, while also claiming their private key remains secret. Further, some non-repudiation schemes offer a timestamp for the digital signature, so that even if the private key is exposed, the signature is valid. Digitally signed messages may be anything representable as a bitstring: examples include electronic mail, contracts, or a message sent via some other cryptographic protocol.
Definition
A digital signature scheme typically consists of three algorithms:
A key generation algorithm that selects a private key uniformly at random from a set of possible private keys. The algorithm outputs the private key and a corresponding public key.
A signing algorithm that, given a message and a private key, produces a signature.
A signature verifying algorithm that, given the message, public key and signature, either accepts or rejects the message's claim to authenticity.
Two main properties are required:
First, the authenticity of a signature generated from a fixed message and fixed private key can be verified by using the corresponding public key.
Secondly, it should be computationally infeasible to generate a valid signature for a party without knowing that party's private key.
A digital signature is an authentication mechanism that enables the creator of the message to attach a code that acts as a signature.
The Digital Signature Algorithm (DSA), developed by the National Institute of Standards and Technology, is one of many examples of a signing algorithm.
In the following discussion, 1n refers to a unary number.
Formally, a digital signature scheme is a triple of probabilistic polynomial time algorithms, (G, S, V), satisfying:
G (key-generator) generates a public key (pk), and a corresponding private key (sk), on input 1n, where n is the security parameter.
S (signing) returns a tag, t, on the inputs: the private key (sk), and a string (x).
V (verifying) outputs accepted or rejected on the inputs: the public key (pk), a string (x), and a tag (t).
For correctness, S and V must satisfy
Pr [ (pk, sk) ← G(1n), V( pk, x, S(sk, x) ) = accepted ] = 1.
A digital signature scheme is secure if for every non-uniform probabilistic polynomial time adversary, A
Pr [ (pk, sk) ← G(1n), (x, t) ← AS(sk, · )(pk, 1n), x ∉ Q, V(pk, x, t) = accepted] < negl(n),
where AS(sk, · ) denotes that A has access to the oracle, S(sk, · ), Q denotes the set of the queries on S made by A, which knows the public key, pk, and the security parameter, n, and x ∉ Q denotes that the adversary may not directly query the string, x, on S.
History
In 1976, Whitfield Diffie and Martin Hellman first described the notion of a digital signature scheme, although they only conjectured that such schemes existed based on functions that are trapdoor one-way permutations. Soon afterwards, Ronald Rivest, Adi Shamir, and Len Adleman invented the RSA algorithm, which could be used to produce primitive digital signatures (although only as a proof-of-concept – "plain" RSA signatures are not secure). The first widely marketed software package to offer digital signature was Lotus | Technology | Computer security | null |
59660 | https://en.wikipedia.org/wiki/Integumentary%20system | Integumentary system | The integumentary system is the set of organs forming the outermost layer of an animal's body. It comprises the skin and its appendages, which act as a physical barrier between the external environment and the internal environment that it serves to protect and maintain the body of the animal. Mainly it is the body's outer skin.
The integumentary system includes skin, hair, scales, feathers, hooves, claws, and nails. It has a variety of additional functions: it may serve to maintain water balance, protect the deeper tissues, excrete wastes, and regulate body temperature, and is the attachment site for sensory receptors which detect pain, sensation, pressure, and temperature.
Structure
Skin
The skin is one of the largest organs of the body. In humans, it accounts for about 12 to 15 percent of total body weight and covers 1.5 to 2 m2 of surface area.
The skin (integument) is a composite organ, made up of at least two major layers of tissue: the epidermis and the dermis. The epidermis is the outermost layer, providing the initial barrier to the external environment. It is separated from the dermis by the basement membrane (basal lamina and reticular lamina). The epidermis contains melanocytes and gives color to the skin. The deepest layer of the epidermis also contains nerve endings. Beneath this, the dermis comprises two sections, the papillary and reticular layers, and contains connective tissues, vessels, glands, follicles, hair roots, sensory nerve endings, and muscular tissue.
Between the integument and the deep body musculature there is a transitional subcutaneous zone made up of very loose connective and adipose tissue, the hypodermis. Substantial collagen bundles anchor the dermis to the hypodermis in a way that permits most areas of the skin to move freely over the deeper tissue layers.
Epidermis
The epidermis is the strong, superficial layer that serves as the first line of protection against the outer environment. The human epidermis is composed of stratified squamous epithelial cells, which further break down into four to five layers: the stratum corneum, stratum granulosum, stratum spinosum and stratum basale. Where the skin is thicker, such as in the palms and soles, there is an extra layer of skin between the stratum corneum and the stratum granulosum, called the stratum lucidum. The epidermis is regenerated from the stem cells found in the basal layer that develop into the corneum. The epidermis itself is devoid of blood supply and draws its nutrition from its underlying dermis.
Its main functions are protection, absorption of nutrients, and homeostasis. In structure, it consists of a keratinized stratified squamous epithelium; four types of cells: keratinocytes, melanocytes, Merkel cells, and Langerhans cells.
The predominant cell keratinocyte, which produces keratin, a fibrous protein that aids in skin protection, is responsible for the formation of the epidermal water barrier by making and secreting lipids. The majority of the skin on the human body is keratinized, with the exception of the lining of mucous membranes, such as the inside of the mouth. Non-keratinized cells allow water to "stay" atop the structure.
The protein keratin stiffens epidermal tissue to form fingernails. Nails grow from a thin area called the nail matrix at an average of 1 mm per week. The lunula is the crescent-shape area at the base of the nail, lighter in color as it mixes with matrix cells. Only primates have nails. In other vertebrates, the keratinizing system at the terminus of each digit produces claws or hooves.
The epidermis of vertebrates is surrounded by two kinds of coverings, which are produced by the epidermis itself. In fish and aquatic amphibians, it is a thin mucus layer that is constantly being replaced. In terrestrial vertebrates, it is the stratum corneum (dead keratinized cells). The epidermis is, to some degree, glandular in all vertebrates, but more so in fish and amphibians. Multicellular epidermal glands penetrate the dermis, where they are surrounded by blood capillaries that provide nutrients and, in the case of endocrine glands, transport their products.
Dermis
The dermis is the underlying connective tissue layer that supports the epidermis. It is composed of dense irregular connective tissue and areolar connective tissue such as a collagen with elastin arranged in a diffusely bundled and woven pattern.
The dermis has two layers: the papillary dermis and the reticular layer. The papillary layer is the superficial layer that forms finger-like projections into the epidermis (dermal papillae), and consists of highly vascularized, loose connective tissue. The reticular layer is the deep layer of the dermis and consists of the dense irregular connective tissue. These layers serve to give elasticity to the integument, allowing stretching and conferring flexibility, while also resisting distortions, wrinkling, and sagging. The dermal layer provides a site for the endings of blood vessels and nerves. Many chromatophores are also stored in this layer, as are the bases of integumental structures such as hair, feathers, and glands.
Hypodermis
The hypodermis, otherwise known as the subcutaneous layer, is a layer beneath the skin. It invaginates into the dermis and is attached to the latter, immediately above it, by collagen and elastin fibers. It is essentially composed of a type of cell known as adipocytes, which are specialized in accumulating and storing fats. These cells are grouped together in lobules separated by connective tissue.
The hypodermis acts as an energy reserve. The fats contained in the adipocytes can be put back into circulation, via the venous route, during intense effort or when there is a lack of energy-providing substances, and are then transformed into energy. The hypodermis participates, passively at least, in thermoregulation since fat is a heat insulator.
Functions
The integumentary system has multiple roles in maintaining the body's equilibrium. All body systems work in an interconnected manner to maintain the internal conditions essential to the function of the body. The skin has an important job of protecting the body and acts as the body's first line of defense against infection, temperature change, and other challenges to homeostasis.
Its main functions include:
Protect the body's internal living tissues and organs
Protect against invasion by infectious organisms
Protect the body from dehydration
Protect the body against abrupt changes in temperature, maintain homeostasis
Help excrete waste materials through perspiration
Act as a receptor for touch, pressure, pain, heat, and cold (see Somatosensory system)
Protect the body against sunburns by secreting melanin
Generate vitamin D through exposure to ultraviolet light
Store water, fat, glucose, vitamin D
Maintenance of the body form
Formation of new cells from stratum germinativum to repair minor injuries
Protect from UV rays.
Regulates body temperature
It distinguishes, separates, and protects the organism from its surroundings.
Small-bodied invertebrates of aquatic or continually moist habitats respire using the outer layer (integument). This gas exchange system, where gases simply diffuse into and out of the interstitial fluid, is called integumentary exchange.
Clinical significance
Possible diseases and injuries to the human integumentary system include:
Rash
Yeast
Athlete's foot
Infection
Sunburn
Skin cancer
Albinism
Acne
Herpes
Herpes labialis, commonly called cold sores
Impetigo
Rubella
Cancer
Psoriasis
Rabies
Rosacea
Atopic dermatitis
Eczema
| Biology and health sciences | Integumentary system | null |
59701 | https://en.wikipedia.org/wiki/Broom | Broom | A broom (also known as a broomstick) is a cleaning tool consisting of usually stiff fibers (often made of materials such as plastic, hair, or corn husks) attached to, and roughly parallel to, a cylindrical handle, the broomstick. It is thus a variety of brush with a long handle. It is commonly used in combination with a dustpan.
A distinction is made between a "hard broom" and a "soft broom" and a spectrum in between. Soft brooms are used in some cultures chiefly for sweeping walls of cobwebs and spiders, like a "feather duster", while hard brooms are for rougher tasks like sweeping dirt off sidewalks or concrete floors, or even smoothing and texturing wet concrete. The majority of brooms are somewhere in between, suitable for sweeping the floors of homes and businesses, soft enough to be flexible and to move even light dust, but stiff enough to achieve a firm sweeping action.
The broom is also a symbolic object associated with witchcraft and ceremonial magic.
Etymology
The word broom derives from types of shrubs referred to as brooms. Common broom typically refers to whatever shrub is most commonly used to make the bristles for a broomstick in a given region. The name of the shrubs began to be used for the household implement in Late Middle English and gradually replaced the earlier besom during the Early Modern English period. The song Buy Broom Buzzems (by William Purvis 1752–1832) still refers to the "broom besom" as one type of besom (i.e. "a besom made from broom").
Flat brooms, made of broom corn, were invented by Shakers in the 19th century with the invention of the broom vice.
A smaller whisk broom or brush is sometimes called a duster.
Function
Brooms are used to clean dust and ash. They may be used to clean homes, appliances such as ovens and fireplaces, or outdoor areas such as streets and yards.
History
The earliest brooms and brushes are from prehistory, when things such as bird wings and burs were fastened to handles of bone, ivory, or wood. The indigenous peoples of the Southwestern United States created brooms from yucca plants for cleaning pueblos. The indigenous people of Saint Lucia created brooms from coconut fronds for cleaning around hearths. Brooms are mentioned in the 1540 manuscript Codex Mendoza of the Aztecs, which instructs girls to sweep.
The birch besom was made by fastening twigs to a handle with a strip of ash wood, harvested from a log after washing it in a running stream. The besom became a symbol of breweries in England, where brewers used it as a whisk while fermenting alcoholic beverages, and the brooms were typically displayed by pubs. When not in use, a brewer's besom was stored and dried on wall pegs or hanging by a leather cord. The broom was not washed so that yeast would remain in the bristles for future uses. Hearth besoms were created in Ireland to keep ash on a hearth. Until the 18th century, brooms were crafted by hand.
In 1797, the quality of brooms changed when Levi Dickenson, a farmer in Hadley, Massachusetts, made a broom for his wife, using the tassels of sorghum, a grain he was growing for the seeds. His wife spread good words around town, creating demand for Dickenson's sorghum brooms. The sorghum brooms held up well, but ultimately, like all brooms, fell apart. Dickenson subsequently invented a machine that would make better brooms, and faster than he could. In 1810, the foot treadle broom machine was invented. This machine played an integral part in the Industrial Revolution.
The Shakers began growing broom corn to create brooms in the present-day United States, which they crafted on treadle wheels and stored hanging on the wall under a cotton hood. The Shaker Theodore Bates invented the flat broom in 1798. Benjamin Franklin grew French broom, a practice which was then taken up by Thomas Jefferson, who had broomsticks made from the plant. Americans commonly kept brooms with their fireplaces by the early 19th century. At this time, brooms were often made by children, the disabled, the elderly, and slaves. By the middle of the century, brooms were created in factories with machine presses, trimmers, and winding machines and then sold door-to-door. People in the American frontier crafted brooms with a wet rawhide fastening, which dried and hardened around the bristles.
Henry Hadley invented a hybridized machine-harvested broom corn at the University of Illinois in 1983 for more efficient creation of brooms. Modern factory-made brooms are made with straw bristles, which are flattened and stitched together before a handle is inserted. In industrialized countries, brooms are sometimes replaced or superseded by powered cleaning instruments such as leaf blowers and vacuum cleaners. Brooms remain commonly used for cleaning purposes in the 21st century.
One source mentions that the United States had 303 broom factories by 1839 and that the number peaked at 1,039 in 1919. Most of these were in the Eastern United States; during the Great Depression in the 1930s, the number of factories declined to 320 in 1939. The state of Oklahoma became a major center for broom production because broom corn grew especially well there, with The Oklahoma Broom Corn Company opening a factory in El Reno in 1906. Faced with competition from imported brooms and synthetic bristles, most of the factories closed by the 1960s.
Design and types
A broom is made up of two parts: the handle, which is a long cylindrical stick, and the stiff fibers lined parallel at its base.
The United States International Cooperation Administration made a distinction between brooms based on bristle quality. Parlor brooms are made of smooth green fibers and typically have brushes 14 to 18 inches long. Carpet brooms are a cheaper variant of the parlor broom that uses bristles rejected for use in parlor brooms for being off-color or lower quality. Standard brooms use bristles that were deemed too low-quality for either parlor brooms or carpet brooms, often dyed green to emulate other brooms. Hearth brooms, or toy brooms, are made of miscellaneous fibers that cannot be used in other brooms. They are not typically sold as consumer products.
Warehouse brooms use heavier fibers such as rattan or palmyra palm and are bound with metal. Different grades of warehouse broom are used to denote the surface it is designed for, such as smelters, decks, or railroads. Their brushes measure about 16 to 18 inches long.
Cob brooms are used to clean webs from high areas and were historically made with round brushes. Whisk brooms use bristles that are shorter and finer than other brooms. Rubber brooms were created in the early 20th century to prevent the debris raised when sweeping with straw brooms.
Materials and production
Brush
The brush of a broom is most commonly made with the fibers of broom corn. Other common plant materials used in brooms include palmyra, rice straw, rice root, piassava, grass, sedge, and twigs. They may use a mix of materials, with lower quality fibers filling out the brush. Broom making involves botanical knowledge, particularly about broom plants.
For manufactured brooms, the fibers are sorted by quality and fitted into the appropriate type of broom. They are then put through an evener to align the fibers, a saw to remove stems, and a scraper to break open the straw and remove the seeds. The fibers are dyed or bleached to achieve a uniform color, or they are wetted if they are already high quality so they can be more easily wound. The outer fibers of the brush are typically treated with a dye, called broom crystals, to preserve the color after use.
As an alternative to plant fibers, brooms can be fitted with synthetic brushes made of materials like nylon or plastic.
Handle and fastening
Wooden broom handles are commonly made from hardwood or fir. Commercial wood broom handles are painted or finished. Lacquers can increase the lifespan of the broom's handle in addition to serving an aesthetic purpose. Wooden broom handles are often about 42 inches long and seven-eighths to one and one-eighth inches in diameter.
Metal tension wires, sometimes crafted specifically for use in brooms, are put through a winding machine to fasten the bristles to the handle. The wire is wound through a hole in the handle before fastening the brush, typically over the last six inches of the handle. Additional bristles are added to the sides for a flat brush shape and to provide a surface for sweeping. The stem ends of the fibers are then cut and tapered and the wire is nailed into the handle. The wire is then finished by one of several methods, such as with a metal cap, with a velvet coat, or by being tapered. After the broom is wired, the fibers can again be scraped or seeded.
Twine, often made of cotton or linen, is used to stitch the brush. At least five stitches will typically be used. The outside of the brush may be wrapped with a material like leather, replacing a twine band used to hold the brush together during manufacturing.
Commercially sold brooms may apply a glued label to the fastening with the brand name or broom model, which can be used as a cover for the clamp marks left by a wiring machine.
Magic
In the context of witchcraft, broomstick is likely to refer to the broom as a whole, known as a besom. The first known reference to witches flying on broomsticks dates to the 11th-century Islamic traditionalist theologian Ibn Qudamahin his book al-Mughnī ( The Persuader ). The first reference to witches flying on broomsticks in Europe dates to 1453, confessed by the male witch Guillaume Edelin. The concept of a flying ointment used by witches appears at about the same time, recorded in 1456.
In Metro-Goldwyn-Mayer's 1939 film, The Wizard of Oz, the Wicked Witch of the West used a broomstick to fly over Oz. She also used it to skywrite "Surrender Dorothy" above the Emerald City. The Wizard commands Dorothy and her three traveling companions to bring the Wicked Witch's broomstick to him in order to grant their wishes. Dorothy carries it to the Wizard with the Scarecrow, Tin Man, and Lion after the Wicked Witch's death.
In Disney's 1940 film Fantasia, Mickey Mouse, playing The Sorcerer's Apprentice, brings a broom to life to do his chore of filling a well full of water. The broom overdoes its job and when chopped into pieces, each splinter becomes a new broom that flood the room until Yen Sid stops them. This story comes from a poem by Goethe called Der Zauberlehrling ("The Sorcerer's Apprentice"). The Disney brooms have had recurring cameos in Disney media, mostly portrayed as janitors, albeit not out of control or causing chaos such as in the original appearance.
This flight was also in Bedknobs and Broomsticks as well as Hocus Pocus.
In Eswatini (Swaziland), witches' broomsticks are short bundles of sticks tied together without a handle.
Flying brooms play an important role in the fantasy world of Harry Potter, used for transportation as well as for playing the popular airborne game of Quidditch. Flying brooms, along with Flying carpets, are the main means of transportation in the world of Poul Anderson's
Operation Chaos.
The Flying Broom () is a feminist organization in Turkey, deliberately evoking the associations of a Flying Broom with witches.
Culture
Brooms are used in some rituals.
Jumping the broom is a tradition sometimes practiced in African American weddings in which the couple leaps over a broom to symbolically represent the leap into domestic life. The tradition was practiced by enslaved Americans and other groups of low social class in the United States through the 19th century. It was revitalized by Alex Haley after it was prominently featured in his novel Roots: The Saga of an American Family in 1976 and became part of a broader reclamation of Black heritage at the time. Other marginalized groups, such as the Celts and the Romani, have historically been described as practicing similar traditions in Britain. The precise origin of jumping the broom is uncertain.
The Métis people of Canada have a broom dancing tradition. There are broom dancing exhibitions where people show off their broom dancing skills. The lively broom dance involves fast footwork and jumping.
During World War II, American submarine crews would tie a broom to their boat's conning tower when returning to port to indicate that they had "swept" the seas clean of enemy shipping. The tradition has been devalued in recent years by submarine crews who fly a broom simply when returning from their boat's shake-down cruise. This tradition may stem from the action of the Dutch admiral Maarten Tromp who tied a broom to his main mast after defeating the British admiral Robert Blake at the Battle of Dungeness in 1652. This has often been interpreted as a message that he would "sweep the British from the seas". This story remains unsubstantiated, but may have its origin in the tradition of hoisting a broom as a sign that a ship was for sale, which seems more likely as Tromp had captured two of Blake's ships in the battle.
In Bhojpuri, it is called Baṛhanī (prosperer), as it is believed that it's prospers the family and house.
Literature
In 1701 Jonathan Swift wrote a "Meditation Upon a Broomstick", a parody of Robert Boyle's Occasional Reflections upon Several Subjects:
In J.K. Rowling's Harry Potter novels and film adaptations, broomsticks are a common form of transport for wizards and witches. These are also used for the magical sport of Quidditch, in which players use their broomsticks to fly around a field and shoot goals.
Politics
For much of the 20th century, political cartoons and propaganda would often depict new or oncoming leaders sweeping away old, corrupt or unpopular figures.
The broom is used as a symbol of the following political parties:
Aam Aadmi Party, India
All Progressives Congress, Nigeria
Religion
In Jainism, monks and nuns have a little broom with them, in order to gently brush aside ants and small animals, to avoid crushing them. This is part of observing the principle of Ahinsā.
The Shakers are often credited with the invention of the flat broom.
Sports
Curling broom
In baseball and basketball, when the home team is close to accomplishing a sweep (having won the first two games of a three-game series or first three games of a four-game series), some fans will bring brooms to the ballpark and brandish them as a way of taunting the visiting team (examples: Arkansas vs. LSU, 2011; Red Sox vs. Yankees, May 13–15, 2011 and June 7–9, 2011).
In broomball, broomsticks have their heads removed and are used to push a ball into a goal, on an ice surface. The game is similar to hockey, except players do not wear skates.
Image gallery
| Technology | Hand tools | null |
59715 | https://en.wikipedia.org/wiki/Scientific%20notation | Scientific notation | Scientific notation is a way of expressing numbers that are too large or too small to be conveniently written in decimal form, since to do so would require writing out an inconveniently long string of digits. It may be referred to as scientific form or standard index form, or standard form in the United Kingdom. This base ten notation is commonly used by scientists, mathematicians, and engineers, in part because it can simplify certain arithmetic operations. On scientific calculators, it is usually known as "SCI" display mode.
In scientific notation, nonzero numbers are written in the form
or m times ten raised to the power of n, where n is an integer, and the coefficient m is a nonzero real number (usually between 1 and 10 in absolute value, and nearly always written as a terminating decimal). The integer n is called the exponent and the real number m is called the significand or mantissa. The term "mantissa" can be ambiguous where logarithms are involved, because it is also the traditional name of the fractional part of the common logarithm. If the number is negative then a minus sign precedes m, as in ordinary decimal notation. In normalized notation, the exponent is chosen so that the absolute value (modulus) of the significand m is at least 1 but less than 10.
Decimal floating point is a computer arithmetic system closely related to scientific notation.
History
Styles
Normalized notation
Any real number can be written in the form in many ways: for example, 350 can be written as or or .
In normalized scientific notation (called "standard form" in the United Kingdom), the exponent n is chosen so that the absolute value of m remains at least one but less than ten (). Thus 350 is written as . This form allows easy comparison of numbers: numbers with bigger exponents are (due to the normalization) larger than those with smaller exponents, and subtraction of exponents gives an estimate of the number of orders of magnitude separating the numbers. It is also the form that is required when using tables of common logarithms. In normalized notation, the exponent n is negative for a number with absolute value between 0 and 1 (e.g. 0.5 is written as ). The 10 and exponent are often omitted when the exponent is 0. For a series of numbers that are to be added or subtracted (or otherwise compared), it can be convenient to use the same value of m for all elements of the series.
Normalized scientific form is the typical form of expression of large numbers in many fields, unless an unnormalized or differently normalized form, such as engineering notation, is desired. Normalized scientific notation is often called exponential notation – although the latter term is more general and also applies when m is not restricted to the range 1 to 10 (as in engineering notation for instance) and to bases other than 10 (for example, ).
Engineering notation
Engineering notation (often named "ENG" on scientific calculators) differs from normalized scientific notation in that the exponent n is restricted to multiples of 3. Consequently, the absolute value of m is in the range 1 ≤ |m| < 1000, rather than 1 ≤ |m| < 10. Though similar in concept, engineering notation is rarely called scientific notation. Engineering notation allows the numbers to explicitly match their corresponding SI prefixes, which facilitates reading and oral communication. For example, can be read as "twelve-point-five nanometres" and written as , while its scientific notation equivalent would likely be read out as "one-point-two-five times ten-to-the-negative-eight metres".
E notation
Calculators and computer programs typically present very large or small numbers using scientific notation, and some can be configured to uniformly present all numbers that way. Because superscript exponents like 107 can be inconvenient to display or type, the letter "E" or "e" (for "exponent") is often used to represent "times ten raised to the power of", so that the notation for a decimal significand m and integer exponent n means the same as . For example is written as or , and is written as or . While common in computer output, this abbreviated version of scientific notation is discouraged for published documents by some style guides.
Most popular programming languages – including Fortran, C/C++, Python, and JavaScript – use this "E" notation, which comes from Fortran and was present in the first version released for the IBM 704 in 1956. The E notation was already used by the developers of SHARE Operating System (SOS) for the IBM 709 in 1958. Later versions of Fortran (at least since FORTRAN IV as of 1961) also use "D" to signify double precision numbers in scientific notation, and newer Fortran compilers use "Q" to signify quadruple precision. The MATLAB programming language supports the use of either "E" or "D".
The ALGOL 60 (1960) programming language uses a subscript ten "10" character instead of the letter "E", for example: 6.0221023. This presented a challenge for computer systems which did not provide such a character, so ALGOL W (1966) replaced the symbol by a single quote, e.g. 6.022'+23, and some Soviet ALGOL variants allowed the use of the Cyrillic letter "ю", e.g. . Subsequently, the ALGOL 68 programming language provided a choice of characters: , , , , or 10. The ALGOL "10" character was included in the Soviet GOST 10859 text encoding (1964), and was added to Unicode 5.2 (2009) as .
Some programming languages use other symbols. For instance, Simula uses (or for long), as in . Mathematica supports the shorthand notation (reserving the letter for the mathematical constant e).
The first pocket calculators supporting scientific notation appeared in 1972. To enter numbers in scientific notation calculators include a button labeled "EXP" or "×10x", among other variants. The displays of pocket calculators of the 1970s did not display an explicit symbol between significand and exponent; instead, one or more digits were left blank (e.g. 6.022 23, as seen in the HP-25), or a pair of smaller and slightly raised digits were reserved for the exponent (e.g. 6.022 23, as seen in the Commodore PR100). In 1976, Hewlett-Packard calculator user Jim Davidson coined the term decapower for the scientific-notation exponent to distinguish it from "normal" exponents, and suggested the letter "D" as a separator between significand and exponent in typewritten numbers (for example, ); these gained some currency in the programmable calculator user community. The letters "E" or "D" were used as a scientific-notation separator by Sharp pocket computers released between 1987 and 1995, "E" used for 10-digit numbers and "D" used for 20-digit double-precision numbers. The Texas Instruments TI-83 and TI-84 series of calculators (1996–present) use a small capital E for the separator.
In 1962, Ronald O. Whitaker of Rowco Engineering Co. proposed a power-of-ten system nomenclature where the exponent would be circled, e.g. 6.022 × 103 would be written as "6.022③".
Significant figures
A significant figure is a digit in a number that adds to its precision. This includes all nonzero numbers, zeroes between significant digits, and zeroes indicated to be significant. Leading and trailing zeroes are not significant digits, because they exist only to show the scale of the number. Unfortunately, this leads to ambiguity. The number is usually read to have five significant figures: 1, 2, 3, 0, and 4, the final two zeroes serving only as placeholders and adding no precision. The same number, however, would be used if the last two digits were also measured precisely and found to equal 0 – seven significant figures.
When a number is converted into normalized scientific notation, it is scaled down to a number between 1 and 10. All of the significant digits remain, but the placeholding zeroes are no longer required. Thus would become if it had five significant digits. If the number were known to six or seven significant figures, it would be shown as or . Thus, an additional advantage of scientific notation is that the number of significant figures is unambiguous.
Estimated final digits
It is customary in scientific measurement to record all the definitely known digits from the measurement and to estimate at least one additional digit if there is any information at all available on its value. The resulting number contains more information than it would without the extra digit, which may be considered a significant digit because it conveys some information leading to greater precision in measurements and in aggregations of measurements (adding them or multiplying them together).
Additional information about precision can be conveyed through additional notation. It is often useful to know how exact the final digit or digits are. For instance, the accepted value of the mass of the proton can properly be expressed as , which is shorthand for . However it is still unclear whether the error ( in this case) is the maximum possible error, standard error, or some other confidence interval.
Use of spaces
In normalized scientific notation, in E notation, and in engineering notation, the space (which in typesetting may be represented by a normal width space or a thin space) that is allowed only before and after "×" or in front of "E" is sometimes omitted, though it is less common to do so before the alphabetical character.
Further examples of scientific notation
An electron's mass is about . In scientific notation, this is written .
The Earth's mass is about . In scientific notation, this is written .
The Earth's circumference is approximately . In scientific notation, this is . In engineering notation, this is written . In SI writing style, this may be written ().
An inch is defined as exactly . Using scientific notation, this value can be uniformly expressed to any desired precision, from the nearest tenth of a millimeter to the nearest nanometer , or beyond.
Hyperinflation means that too much money is put into circulation, perhaps by printing banknotes, chasing too few goods. It is sometimes defined as inflation of 50% or more in a single month. In such conditions, money rapidly loses its value. Some countries have had events of inflation of 1 million percent or more in a single month, which usually results in the rapid abandonment of the currency. For example, in November 2008 the monthly inflation rate of the Zimbabwean dollar reached 79.6 billion percent (470% per day); the approximate value with three significant figures would be %, or more simply a rate of .
Converting numbers
Converting a number in these cases means to either convert the number into scientific notation form, convert it back into decimal form or to change the exponent part of the equation. None of these alter the actual number, only how it's expressed.
Decimal to scientific
First, move the decimal separator point sufficient places, n, to put the number's value within a desired range, between 1 and 10 for normalized notation. If the decimal was moved to the left, append × 10n; to the right, × 10−n. To represent the number in normalized scientific notation, the decimal separator would be moved 6 digits to the left and × 106 appended, resulting in . The number would have its decimal separator shifted 3 digits to the right instead of the left and yield as a result.
Scientific to decimal
Converting a number from scientific notation to decimal notation, first remove the × 10n on the end, then shift the decimal separator n digits to the right (positive n) or left (negative n). The number would have its decimal separator shifted 6 digits to the right and become , while would have its decimal separator moved 3 digits to the left and be .
Exponential
Conversion between different scientific notation representations of the same number with different exponential values is achieved by performing opposite operations of multiplication or division by a power of ten on the significand and an subtraction or addition of one on the exponent part. The decimal separator in the significand is shifted x places to the left (or right) and x is added to (or subtracted from) the exponent, as shown below.
Basic operations
Given two numbers in scientific notation,
and
Multiplication and division are performed using the rules for operation with exponentiation:
and
Some examples are:
and
Addition and subtraction require the numbers to be represented using the same exponential part, so that the significand can be simply added or subtracted:
Next, add or subtract the significands:
An example:
Other bases
While base ten is normally used for scientific notation, powers of other bases can be used too, base 2 being the next most commonly used one.
For example, in base-2 scientific notation, the number 1001b in binary (=9d) is written as or using binary numbers (or shorter if binary context is obvious). In E notation, this is written as (or shorter: 1.001E11) with the letter "E" now standing for "times two (10b) to the power" here. In order to better distinguish this base-2 exponent from a base-10 exponent, a base-2 exponent is sometimes also indicated by using the letter "B" instead of "E", a shorthand notation originally proposed by Bruce Alan Martin of Brookhaven National Laboratory in 1968, as in (or shorter: 1.001B11). For comparison, the same number in decimal representation: (using decimal representation), or 1.125B3 (still using decimal representation). Some calculators use a mixed representation for binary floating point numbers, where the exponent is displayed as decimal number even in binary mode, so the above becomes or shorter 1.001B3.
This is closely related to the base-2 floating-point representation commonly used in computer arithmetic, and the usage of IEC binary prefixes (e.g. 1B10 for 1×210 (kibi), 1B20 for 1×220 (mebi), 1B30 for 1×230 (gibi), 1B40 for 1×240 (tebi)).
Similar to "B" (or "b"), the letters "H" (or "h") and "O" (or "o", or "C") are sometimes also used to indicate times 16 or 8 to the power as in 1.25 = = 1.40H0 = 1.40h0, or 98000 = = 2.7732o5 = 2.7732C5.
Another similar convention to denote base-2 exponents is using a letter "P" (or "p", for "power"). In this notation the significand is always meant to be hexadecimal, whereas the exponent is always meant to be decimal. This notation can be produced by implementations of the printf family of functions following the C99 specification and (Single Unix Specification) IEEE Std 1003.1 POSIX standard, when using the %a or %A conversion specifiers. Starting with C++11, C++ I/O functions could parse and print the P notation as well. Meanwhile, the notation has been fully adopted by the language standard since C++17. Apple's Swift supports it as well. It is also required by the IEEE 754-2008 binary floating-point standard. Example: 1.3DEp42 represents .
Engineering notation can be viewed as a base-1000 scientific notation.
| Mathematics | Basics | null |
59718 | https://en.wikipedia.org/wiki/Identity%20matrix | Identity matrix | In linear algebra, the identity matrix of size is the square matrix with ones on the main diagonal and zeros elsewhere. It has unique properties, for example when the identity matrix represents a geometric transformation, the object remains unchanged by the transformation. In other contexts, it is analogous to multiplying by the number 1.
Terminology and notation
The identity matrix is often denoted by , or simply by if the size is immaterial or can be trivially determined by the context.
The term unit matrix has also been widely used, but the term identity matrix is now standard. The term unit matrix is ambiguous, because it is also used for a matrix of ones and for any unit of the ring of all matrices.
In some fields, such as group theory or quantum mechanics, the identity matrix is sometimes denoted by a boldface one, , or called "id" (short for identity). Less frequently, some mathematics books use or to represent the identity matrix, standing for "unit matrix" and the German word respectively.
In terms of a notation that is sometimes used to concisely describe diagonal matrices, the identity matrix can be written as
The identity matrix can also be written using the Kronecker delta notation:
Properties
When is an matrix, it is a property of matrix multiplication that
In particular, the identity matrix serves as the multiplicative identity of the matrix ring of all matrices, and as the identity element of the general linear group , which consists of all invertible matrices under the matrix multiplication operation. In particular, the identity matrix is invertible. It is an involutory matrix, equal to its own inverse. In this group, two square matrices have the identity matrix as their product exactly when they are the inverses of each other.
When matrices are used to represent linear transformations from an -dimensional vector space to itself, the identity matrix represents the identity function, for whatever basis was used in this representation.
The th column of an identity matrix is the unit vector , a vector whose th entry is 1 and 0 elsewhere. The determinant of the identity matrix is 1, and its trace is .
The identity matrix is the only idempotent matrix with non-zero determinant. That is, it is the only matrix such that:
When multiplied by itself, the result is itself
All of its rows and columns are linearly independent.
The principal square root of an identity matrix is itself, and this is its only positive-definite square root. However, every identity matrix with at least two rows and columns has an infinitude of symmetric square roots.
The rank of an identity matrix equals the size , i.e.:
| Mathematics | Linear algebra | null |
59733 | https://en.wikipedia.org/wiki/Hexagon | Hexagon | In geometry, a hexagon (from Greek , , meaning "six", and , , meaning "corner, angle") is a six-sided polygon. The total of the internal angles of any simple (non-self-intersecting) hexagon is 720°.
Regular hexagon
A regular hexagon has Schläfli symbol {6} and can also be constructed as a truncated equilateral triangle, t{3}, which alternates two types of edges.
A regular hexagon is defined as a hexagon that is both equilateral and equiangular. It is bicentric, meaning that it is both cyclic (has a circumscribed circle) and tangential (has an inscribed circle).
The common length of the sides equals the radius of the circumscribed circle or circumcircle, which equals times the apothem (radius of the inscribed circle). All internal angles are 120 degrees. A regular hexagon has six rotational symmetries (rotational symmetry of order six) and six reflection symmetries (six lines of symmetry), making up the dihedral group D6. The longest diagonals of a regular hexagon, connecting diametrically opposite vertices, are twice the length of one side. From this it can be seen that a triangle with a vertex at the center of the regular hexagon and sharing one side with the hexagon is equilateral, and that the regular hexagon can be partitioned into six equilateral triangles.
Like squares and equilateral triangles, regular hexagons fit together without any gaps to tile the plane (three hexagons meeting at every vertex), and so are useful for constructing tessellations. The cells of a beehive honeycomb are hexagonal for this reason and because the shape makes efficient use of space and building materials. The Voronoi diagram of a regular triangular lattice is the honeycomb tessellation of hexagons.
Parameters
The maximal diameter (which corresponds to the long diagonal of the hexagon), D, is twice the maximal radius or circumradius, R, which equals the side length, t. The minimal diameter or the diameter of the inscribed circle (separation of parallel sides, flat-to-flat distance, short diagonal or height when resting on a flat base), d, is twice the minimal radius or inradius, r. The maxima and minima are related by the same factor:
and, similarly,
The area of a regular hexagon
For any regular polygon, the area can also be expressed in terms of the apothem a and the perimeter p. For the regular hexagon these are given by a = r, and p, so
The regular hexagon fills the fraction of its circumscribed circle.
If a regular hexagon has successive vertices A, B, C, D, E, F and if P is any point on the circumcircle between B and C, then .
It follows from the ratio of circumradius to inradius that the height-to-width ratio of a regular hexagon is 1:1.1547005; that is, a hexagon with a long diagonal of 1.0000000 will have a distance of 0.8660254 or cos(30°) between parallel sides.
Point in plane
For an arbitrary point in the plane of a regular hexagon with circumradius , whose distances to the centroid of the regular hexagon and its six vertices are and
respectively, we have
If are the distances from the vertices of a regular hexagon to any point on its circumcircle, then
Symmetry
The regular hexagon has D6 symmetry. There are 16 subgroups. There are 8 up to isomorphism: itself (D6), 2 dihedral: (D3, D2), 4 cyclic: (Z6, Z3, Z2, Z1) and the trivial (e)
These symmetries express nine distinct symmetries of a regular hexagon. John Conway labels these by a letter and group order. r12 is full symmetry, and a1 is no symmetry. p6, an isogonal hexagon constructed by three mirrors can alternate long and short edges, and d6, an isotoxal hexagon constructed with equal edge lengths, but vertices alternating two different internal angles. These two forms are duals of each other and have half the symmetry order of the regular hexagon. The i4 forms are regular hexagons flattened or stretched along one symmetry direction. It can be seen as an elongated rhombus, while d2 and p2 can be seen as horizontally and vertically elongated kites. g2 hexagons, with opposite sides parallel are also called hexagonal parallelogons.
Each subgroup symmetry allows one or more degrees of freedom for irregular forms. Only the g6 subgroup has no degrees of freedom but can be seen as directed edges.
Hexagons of symmetry g2, i4, and r12, as parallelogons can tessellate the Euclidean plane by translation. Other hexagon shapes can tile the plane with different orientations.
A2 and G2 groups
The 6 roots of the simple Lie group A2, represented by a Dynkin diagram , are in a regular hexagonal pattern. The two simple roots have a 120° angle between them.
The 12 roots of the Exceptional Lie group G2, represented by a Dynkin diagram are also in a hexagonal pattern. The two simple roots of two lengths have a 150° angle between them.
Dissection
Coxeter states that every zonogon (a 2m-gon whose opposite sides are parallel and of equal length) can be dissected into parallelograms. In particular this is true for regular polygons with evenly many sides, in which case the parallelograms are all rhombi. This decomposition of a regular hexagon is based on a Petrie polygon projection of a cube, with 3 of 6 square faces. Other parallelogons and projective directions of the cube are dissected within rectangular cuboids.
Related polygons and tilings
A regular hexagon has Schläfli symbol {6}. A regular hexagon is a part of the regular hexagonal tiling, {6,3}, with three hexagonal faces around each vertex.
A regular hexagon can also be created as a truncated equilateral triangle, with Schläfli symbol t{3}. Seen with two types (colors) of edges, this form only has D3 symmetry.
A truncated hexagon, t{6}, is a dodecagon, {12}, alternating two types (colors) of edges. An alternated hexagon, h{6}, is an equilateral triangle, {3}. A regular hexagon can be stellated with equilateral triangles on its edges, creating a hexagram. A regular hexagon can be dissected into six equilateral triangles by adding a center point. This pattern repeats within the regular triangular tiling.
A regular hexagon can be extended into a regular dodecagon by adding alternating squares and equilateral triangles around it. This pattern repeats within the rhombitrihexagonal tiling.
Self-crossing hexagons
There are six self-crossing hexagons with the vertex arrangement of the regular hexagon:
Hexagonal structures
From bees' honeycombs to the Giant's Causeway, hexagonal patterns are prevalent in nature due to their efficiency. In a hexagonal grid each line is as short as it can possibly be if a large area is to be filled with the fewest hexagons. This means that honeycombs require less wax to construct and gain much strength under compression.
Irregular hexagons with parallel opposite edges are called parallelogons and can also tile the plane by translation. In three dimensions, hexagonal prisms with parallel opposite faces are called parallelohedrons and these can tessellate 3-space by translation.
Tesselations by hexagons
In addition to the regular hexagon, which determines a unique tessellation of the plane, any irregular hexagon which satisfies the Conway criterion will tile the plane.
Hexagon inscribed in a conic section
Pascal's theorem (also known as the "Hexagrammum Mysticum Theorem") states that if an arbitrary hexagon is inscribed in any conic section, and pairs of opposite sides are extended until they meet, the three intersection points will lie on a straight line, the "Pascal line" of that configuration.
Cyclic hexagon
The Lemoine hexagon is a cyclic hexagon (one inscribed in a circle) with vertices given by the six intersections of the edges of a triangle and the three lines that are parallel to the edges that pass through its symmedian point.
If the successive sides of a cyclic hexagon are a, b, c, d, e, f, then the three main diagonals intersect in a single point if and only if .
If, for each side of a cyclic hexagon, the adjacent sides are extended to their intersection, forming a triangle exterior to the given side, then the segments connecting the circumcenters of opposite triangles are concurrent.
If a hexagon has vertices on the circumcircle of an acute triangle at the six points (including three triangle vertices) where the extended altitudes of the triangle meet the circumcircle, then the area of the hexagon is twice the area of the triangle.
Hexagon tangential to a conic section
Let ABCDEF be a hexagon formed by six tangent lines of a conic section. Then Brianchon's theorem states that the three main diagonals AD, BE, and CF intersect at a single point.
In a hexagon that is tangential to a circle and that has consecutive sides a, b, c, d, e, and f,
Equilateral triangles on the sides of an arbitrary hexagon
If an equilateral triangle is constructed externally on each side of any hexagon, then the midpoints of the segments connecting the centroids of opposite triangles form another equilateral triangle.
Skew hexagon
A skew hexagon is a skew polygon with six vertices and edges but not existing on the same plane. The interior of such a hexagon is not generally defined. A skew zig-zag hexagon has vertices alternating between two parallel planes.
A regular skew hexagon is vertex-transitive with equal edge lengths. In three dimensions it will be a zig-zag skew hexagon and can be seen in the vertices and side edges of a triangular antiprism with the same D3d, [2+,6] symmetry, order 12.
The cube and octahedron (same as triangular antiprism) have regular skew hexagons as petrie polygons.
Petrie polygons
The regular skew hexagon is the Petrie polygon for these higher dimensional regular, uniform and dual polyhedra and polytopes, shown in these skew orthogonal projections:
Convex equilateral hexagon
A principal diagonal of a hexagon is a diagonal which divides the hexagon into quadrilaterals. In any convex equilateral hexagon (one with all sides equal) with common side a, there exists a principal diagonal d1 such that
and a principal diagonal d2 such that
Polyhedra with hexagons
There is no Platonic solid made of only regular hexagons, because the hexagons tessellate, not allowing the result to "fold up". The Archimedean solids with some hexagonal faces are the truncated tetrahedron, truncated octahedron, truncated icosahedron (of soccer ball and fullerene fame), truncated cuboctahedron and the truncated icosidodecahedron. These hexagons can be considered truncated triangles, with Coxeter diagrams of the form and .
There are other symmetry polyhedra with stretched or flattened hexagons, like these Goldberg polyhedron G(2,0):
There are also 9 Johnson solids with regular hexagons:
Hexagon versus Sexagon
The debate over whether hexagons should be referred to as "sexagons" has its roots in the etymology of the term. The prefix "hex-" originates from the Greek word "hex," meaning six, while "sex-" comes from the Latin "sex," also signifying six. Some linguists and mathematicians argue that since many English mathematical terms derive from Latin, the use of "sexagon" would align with this tradition. Historical discussions date back to the 19th century, when mathematicians began to standardize terminology in geometry. However, the term "hexagon" has prevailed in common usage and academic literature, solidifying its place in mathematical terminology despite the historical argument for "sexagon." The consensus remains that "hexagon" is the appropriate term, reflecting its Greek origins and established usage in mathematics. (see Numeral_prefix#Occurrences).
Gallery of natural and artificial hexagons
| Mathematics | Two-dimensional space | null |
59746 | https://en.wikipedia.org/wiki/Brown%20recluse%20spider | Brown recluse spider | The brown recluse (Loxosceles reclusa, Sicariidae, formerly placed in a family "Loxoscelidae") is a recluse spider with necrotic venom. Similar to those of other recluse spiders, their bites sometimes require medical attention. The brown recluse is one of three spiders in North America with dangerous venom, the others being the black widow and the Chilean recluse.
Brown recluse spiders are usually between , but may grow larger. While typically light to medium brown, they range in color from whitish to dark brown or blackish gray. The cephalothorax and abdomen are not necessarily the same color. These spiders usually have markings on the dorsal side of their cephalothorax, with a black line coming from it that looks like a violin with the neck of the violin pointing to the rear of the spider, resulting in the nicknames fiddleback spider, brown fiddler, or violin spider.
Description
The violin pattern is not a definitive identifier, as other spiders can have similar markings (e.g. cellar spiders and pirate spiders). Instead, while most spiders have eight eyes, recluse spiders have six eyes arranged in pairs (dyads) with one median pair and two lateral pairs. Only a few other spiders have three pairs of eyes arranged in this way (e.g., scytodids). Recluses have no obvious coloration patterns on the abdomen or legs, and the legs lack spines.
The violin marking can vary in intensity depending on the age of the brown recluse spider, with mature spiders typically having dark violin shapes.
Distribution
The documented range of this species lies roughly south of a line from southeastern Nebraska through southern Iowa, Illinois, and Indiana to southwestern Ohio. In the southern states, it is native from central Texas to western Georgia and north to Kentucky.
Despite rumors to the contrary, the brown recluse spider has not established itself in California or anywhere outside its native range. There are other species of the genus Loxosceles native to the southwestern part of the United States, including California, which may resemble the brown recluse, but interactions between humans and the recluse species in California and the region are rare because those species native ranges lie outside of dense human populations. The number of "false positive" reports based on misidentifications is considerable; in a nationwide study where people submitted spiders that they thought were brown recluses, of 581 from California only one was a brown recluse—submitted by a family that moved from Missouri and brought it with them (compared to specimens submitted from Missouri, Kansas, and Oklahoma, where between 75% and 90% were recluses).
From this study, the most common spider submitted from California as a brown recluse was in the genus Titiotus, whose bite is deemed harmless. A similar study documented that various arachnids were routinely misidentified by physicians, pest control operators, and other non-expert authorities, who told their patients or clients that the spider they had was a brown recluse when in fact it was not.
Despite the absence of brown recluses from the Western U.S., physicians in the region commonly diagnose "brown recluse bites", leading to the popular misconception that the spiders inhabit those areas.
Over the last century, spiders have occasionally been intercepted in locations where they have no known established populations; these spiders may be transported fairly easily, though the lack of established populations well outside the natural range also indicates that such movement has not led to the colonization of new areas, after decades of opportunities. Note that the occurrence of brown recluses in a single building (such as a warehouse) outside of the native range is not considered as successful colonization; such single-building populations can occur (e.g., in several such cases in Florida),
but do not spread, and can be easily eradicated.
The spider has also received numerous sensationalized media reports of bites occurring where these spiders are absent (and no specimens were found), such as a 2014 report from Thailand, where a man was claimed to have died from a brown recluse bite.
Many misidentifications and erroneous geographic records stem from the similarity between L. reclusa and a related introduced species, the Mediterranean recluse (Loxosceles rufescens), which is found worldwide, including numerous sightings throughout the United States; the two species are superficially almost indistinguishable, and misidentifications are common, making it difficult to distinguish which reports of recluses refer to which species.
Life cycle
Adult brown recluse spiders often live about one to two years. Each female produces several egg sacs over a period of two to three months, from May to July, with approximately 50 eggs in each sac. The eggs hatch in about one month. The spiderlings take about one year to grow to adulthood. The brown recluse spider is resilient and can tolerate up to six months of extreme drought and scarcity or absence of food. On one occasion, a brown recluse survived in controlled captivity for over five seasons without any food at all.
As part of the haplogynae, brown recluses do not balloon, which limits their ability to widely disperse geographically.
The brown recluse will, though not habitually, cannibalize another if food becomes scarce; especially during its typical mating season from June to September or when an unreceptive female encounters an aggressive male.
Behavior
A brown recluse's stance on a flat surface is usually with all legs radially extended. When alarmed it may lower its body, withdraw the forward two legs straight rearward into a defensive position, withdraw the rearmost pair of legs into a position for lunging forward, and stand motionless with pedipalps raised. The pedipalps in mature specimens are dark and quite prominent and are normally held horizontally forward. When threatened it usually flees, seemingly to avoid a conflict, and if detained may further avoid contact with quick horizontal rotating movements or even resort to assuming a lifeless pose (playing dead). The spider does not usually jump unless touched brusquely, and even then its avoidance movement is more of a horizontal lunge rather than a vaulting of itself entirely off the surface. When running, the brown recluse does not leave a silk line behind, which would make it more easily tracked when it is being pursued. Movement at virtually any speed is an evenly paced gait with legs extended.
The brown recluse spider displays autotomy as a defense mechanism against physical, predatory attack to a leg as well as to prevent predatory, venom injections from spreading to the rest of the body. "Once a leg is lost, a recluse spider does not regenerate a new one with subsequent molts", unlike the huntsman spider, which does regenerate autotomized legs. With each time that a leg is autotomized, the recluse "changes its gait to compensate for the loss."
Habitat
Brown recluse spiders build asymmetrical (irregular) webs that frequently include a shelter consisting of disorderly threads. They frequently build their webs in woodpiles and sheds, closets, garages, plenum spaces, cellars, and other places that are dry and generally undisturbed. When dwelling in human residences they seem to favor cardboard, possibly because it mimics the rotting tree bark which they inhabit naturally. Human-recluse contact often occurs when such isolated spaces are disturbed and the spider feels threatened. Unlike most web weavers, they leave these lairs at night to hunt. Males move around more when hunting than the females, which tend to remain nearer to their webs.
Bite
Like all members of the Loxosceles genus, the brown recluse has potent tissue-destroying venoms containing the dermonecrotic agent sphingomyelinase D.
Most bites are minor with no dermonecrosis, but a small number of brown recluse bites produce loxoscelism, a condition where the skin around the bite dies. While loxoscelism usually manifests as a skin condition (cutaneous loxoscelism), it can also include systemic symptoms like fever, nausea, and vomiting (viscerocutaneous loxoscelism). In very rare cases, bites can even cause hemolysis – the bursting of red blood cells.
In one study of clinically diagnosed brown recluse bites, skin necrosis occurred 37% of the time, while systemic illness occurred 14% of the time.
As suggested by its specific name reclusa (recluse), the brown recluse spider is rarely aggressive, and bites from the species are uncommon. In 2001, more than 2,000 brown recluse spiders were removed from a heavily infested home in Kansas, yet the four residents who had lived there for years were never harmed by the spiders, despite many encounters with them.
The spider usually bites only when pressed against the skin, such as when tangled within clothes, shoes, towels, bedding, inside work gloves, etc. Many human victims report having been bitten after putting on clothes or shoes that had not been worn recently or had been left for many days undisturbed on the floor. The fangs of the brown recluse are not large enough to penetrate most fabric.
When both types of loxoscelism do result, systemic effects may occur before necrosis, as the venom spreads throughout the body in minutes. Children, the elderly, and the debilitatingly ill may be more susceptible to systemic loxoscelism. The systemic symptoms most commonly experienced include nausea, vomiting, fever, rashes, and muscle and joint pain. Rarely, such bites can result in hemolysis, low platelet levels, blood clots throughout the body, organ damage, and even death.
Most fatalities are in children under the age of seven
or those with a weak immune system.
While the majority of brown recluse spider bites do not result in any symptoms, cutaneous symptoms occur more frequently than systemic symptoms. In such instances, the bite forms a necrotizing ulcer as the result of soft tissue destruction and may take months to heal, leaving deep scars. These bites usually become painful and itchy within 2–8 hours. Pain and other local effects worsen 12–36 hours after the bite, and the necrosis develops over the next few days.
Over time, the wound may grow to as large as 25 cm (10 inches). The damaged tissue becomes gangrenous and eventually sloughs away. L. reclusa can produce slightly more than 0.1 μL of venom, though the average yield is less.
Misdiagnosis
There is an ELISA-based test for brown recluse venom that can determine whether a wound is a brown recluse bite, although it is not commercially available and not in routine clinical use. Clinical diagnoses often use Occam's razor principle in diagnosing bites based on what spiders the patient likely encountered and previous similar diagnoses.
suggested the mnemonic "NOT RECLUSE", shown below, as a memory device to help laymen and medical professionals more objectively screen and diagnose potential cases of loxoscelism.
{| style="text-align:left;vertical-align:top;"
|+ NOT RECLUSE mnemonic
|-
|colspan=2 align=center|Note that these are contrary criteria: Any one being true indicates that the injury is not a brown recluse bite.
|- style="vertical-align:top;"
| umerous
| More than one wound.
|- style="vertical-align:top;"
| ccurrence
| The injury did not occur in a place where brown recluses are likely to be found:Either outside of the spider's geographic territory ... or not in an enclosed spacelike a box, closet, or attic.
|- style="vertical-align:top;"
| iming
| The wound arose sometime between November and March.
|- style="vertical-align:top;"
| ed center
| The center of the wound is red.
|- style="vertical-align:top;"
| levated
| The middle of the wound is elevated, not sunken.
|- style="vertical-align:top;"
| hronic
| The wound has persisted for more than three months.
|- style="vertical-align:top;"
| arge
| The wound is more than wide.
|- style="vertical-align:top;"
| lcerates too early
| The wound gets crusty within the first week.
|- style="vertical-align:top;"
| wollen
| The wound swells up if it's below the neck.
|- style="vertical-align:top;"
| xudative
| The wound is "wet" – oozing pus or clear fluid.
|}
There are numerous documented infectious and noninfectious conditions that produce wounds that have been initially misdiagnosed as recluse spider bites by medical professionals, including:
Pyoderma gangrenosum
Infection by Staphylococcus
Infection by Streptococcus
Herpes
Diabetic ulcers
Fungal infection
Chemical burns
Toxicodendron dermatitis
Squamous cell carcinoma
Localized vasculitis
Syphilis
Toxic epidermal necrolysis
Sporotrichosis
Lyme disease
Many of these conditions are far more common and more likely to be the source of necrotic wounds, even in areas where brown recluse spiders actually occur. The most important of these is methicillin-resistant Staphylococcus aureus (MRSA), a bacterium whose necrotic lesions are very similar to those induced by recluse bites, and which can be lethal if left untreated. Misdiagnosis of MRSA as spider bites is extremely common (nearly 30% of patients with MRSA reported that they initially suspected a spider bite), and can have fatal consequences.
Reported cases of brown recluse bites occur primarily in Arkansas, Colorado, Kansas, Missouri, Nebraska, Oklahoma, and Texas.
There have been many reports of brown recluse bites in California – though a few related species of spiders may be found there, none of the related spiders in California is known to bite humans.
To date, the reports of bites from areas outside of the spider's native range have been either unverified or, if verified, the spiders have been moved to those locations by travelers or commerce.
Other spiders
For a comparison of the toxicity of several kinds of spider
bites, see the list of medically significant spider bites
Many arachnologists believe that a large number of bites attributed to the brown recluse in the U.S. West Coast are either from other spider species or not spider bites at all. Other spiders in western states that might possibly cause necrotic injuries are the hobo spider, desert recluse spider, and the yellow sac spider.
For example, the venom of the hobo spider, a common European species now established in the northwestern United States and southern British Columbia, has been reported to produce similar symptoms as the brown recluse bite when injected into laboratory rabbits. However, the toxicity of hobo spider venom has been called into question: Actual bites (rather than syringe injections) have not been shown to cause necrosis, and no necrotic hobo spider bites have ever been reported where it is native.
Numerous other spiders have been associated with necrotic bites in the medical literature. Other recluse species, such as the desert recluse (found in the deserts of southwestern United States), are reported to have caused necrotic bite wounds, though only rarely.
The hobo spider and the yellow sac spider have also been reported to cause necrotic bites. However, the bites from these spiders are not known to produce the severe symptoms that can follow from a recluse spider bite, and the level of danger posed by these has been called into question.
So far, no known necrotoxins have been isolated from the venom of any of these spiders, and some arachnologists have disputed the accuracy of spider identifications carried out by bite victims, family members, medical responders, and other non-experts in arachnology.
There have been several studies questioning the danger posed by some of the other spiders mentioned: In these studies, scientists examined case studies of bites in which the spider in question was identified by an expert, and found that the incidence of necrotic injury diminished significantly when "questionable" identifications were excluded from the sample set.
Bite treatment
First aid involves the application of an ice pack to control inflammation and prompt medical care. If it can be easily captured, the spider should be brought with the patient in a clear, tightly closed container so it may be identified by an arachnologist; if there is no specimen at all, then confirmation by an expert is impossible.
Routine treatment should include immobilization of the affected limb, application of ice, local wound care, and tetanus prophylaxis. Many other therapies have been used with varying degrees of success, including hyperbaric oxygen, dapsone, antihistamines (e.g., cyproheptadine), antibiotics, dextran, glucocorticoids, vasodilators, heparin, nitroglycerin, electric shock, curettage, surgical excision, and antivenom.
In almost all cases, bites are self-limited and typically heal without any medical intervention.
Outpatient palliative care following discharge often consists of a weak or moderate strength opioid (e.g. codeine or tramadol, respectively) depending on pain scores, an anti-inflammatory agent (e.g. naproxen, cortisone), and an antispasmodic (e.g. cyclobenzaprine, diazepam), for a few days to a week. If the pain or spasms have not resolved by this time, a second medical evaluation is generally advised, and differential diagnoses may be considered.
Specific treatments
In presumed cases of recluse bites, dapsone is often used for the treatment of necrosis, but controlled clinical trials have yet to demonstrate efficacy.
However, dapsone may be effective in treating many "spider bites" because many such cases are actually misdiagnosed microbial infections.
There have been conflicting reports about its efficacy in treating brown recluse bites, and some have suggested it should no longer be used routinely, if at all.
Wound infection is rare. Antibiotics are not recommended unless there is a credible diagnosis of infection.
Studies have shown that surgical intervention is ineffective and may worsen outcomes. Excision may delay wound healing, cause abscesses, and lead to scarring.
Purportedly application of nitroglycerin stopped necrosis. However, one scientific animal study found no benefit in preventing necrosis, with the study's results showing it increased inflammation and caused symptoms of systemic envenoming. The authors concluded the results of the study did not support the use of topical nitroglycerin in brown recluse envenoming.
Antivenom is available in South America for the venom of related species of recluse spiders. However, the bites, often being painless, usually do not present symptoms until 24 or more hours after the event, possibly limiting the effect of this intervention.
Spider population control
Due to increased fear of these spiders prompted by greater public awareness of their presence in recent years, the extermination of domestic brown recluses is performed frequently in the lower midwestern United States. Brown recluse spiders possess a variety of adaptive abilities, including the ability to survive up to 10 months with no food or water. Additionally, these spiders survive significantly longer in a relatively cool, thermally stable environment.
| Biology and health sciences | Spiders | Animals |
59790 | https://en.wikipedia.org/wiki/Farm | Farm | A farm (also called an agricultural holding) is an area of land that is devoted primarily to agricultural processes with the primary objective of producing food and other crops; it is the basic facility in food production. The name is used for specialized units such as arable farms, vegetable farms, fruit farms, dairy, pig and poultry farms, and land used for the production of natural fiber, biofuel, and other commodities. It includes ranches, feedlots, orchards, plantations and estates, smallholdings, and hobby farms, and includes the farmhouse and agricultural buildings as well as the land. In modern times, the term has been extended to include such industrial operations as wind farms and fish farms, both of which can operate on land or at sea.
There are about 570 million farms in the world, most of which are small and family-operated. Small farms with a land area of fewer than 2 hectares operate on about 12% of the world's agricultural land, and family farms comprise about 75% of the world's agricultural land.
Modern farms in developed countries are highly mechanized. In the United States, livestock may be raised on rangeland and finished in feedlots, and the mechanization of crop production has brought about a great decrease in the number of agricultural workers needed. In Europe, traditional family farms are giving way to larger production units. In Australia, some farms are very large because the land is unable to support a high stocking density of livestock because of climatic conditions. In less developed countries, small farms are the norm, and the majority of rural residents are subsistence farmers, feeding their families and selling any surplus products in the local market.
Etymology
The word in the sense of an agricultural land-holding derives from the verb "to farm" a revenue source, whether taxes, customs, rents of a group of manors or simply to hold an individual manor by the feudal land tenure of "fee farm". The word is from the medieval Latin noun firma, also the source of the French word ferme, meaning a fixed agreement, contract, from the classical Latin adjective firmus meaning strong, stout, firm. As in the medieval age virtually all manors were engaged in the business of agriculture, which was their principal revenue source, so to hold a manor by the tenure of "fee farm" became synonymous with the practice of agriculture itself.
History
Farming has been innovated at multiple different points and places in human history. The transition from hunter-gatherer to settled, agricultural societies is called the Neolithic Revolution and first began around 12,000 years ago, near the beginning of the geological epoch of the Holocene around 12,000 years ago. It was the world's first historically verifiable revolution in agriculture. Farming spread from the Middle East to Europe and by 4,000 BC people that lived in the central part of Europe were using oxen to pull plows and wagons. Subsequent step-changes in human farming practices were provoked by the British Agricultural Revolution in the 18th century, and the Green Revolution of the second half of the 20th century.
Farming originated independently in different parts of the world, as hunter-gatherer societies transitioned to food production rather than food capture. It may have started about 12,000 years ago with the domestication of livestock in the Fertile Crescent in western Asia, soon to be followed by the cultivation of crops. Modern units tend to specialize in the crops or livestock best suited to the region, with their finished products being sold for the retail market or for further processing, with farm products being traded around the world.
Types of farms
A farm may be owned and operated by a single individual, family, community, corporation, or a company, may produce one or many types of produce, and can be a holding of any size from a fraction of a hectare to several thousand hectares.
A farm may operate under a monoculture system or with a variety of cereal or arable crops, which may be separate from or combined with raising livestock. Specialist farms are often denoted as such, thus a dairy farm, fish farm, poultry farm or mink farm.
Some farms may not use the word at all, hence vineyard (grapes), orchard (nuts and other fruit), market garden or "truck farm" (vegetables and flowers). Some farms may be denoted by their topographical location, such as a hill farm, while large estates growing cash crops such as cotton or coffee may be called plantations.
Many other terms are used to describe farms to denote their methods of production, as in collective, corporate, intensive, organic or vertical.
Where most of the income is from some other employment, and the farm is an expanded residence, the term hobby farm is common. This will allow sufficient size for recreational use but be very unlikely to produce sufficient income to be self-sustaining. Hobby farms are commonly around but may be much larger depending on land prices.
Other farms may primarily exist for research or education, such as an ant farm, and since farming is synonymous with mass production, the word "farm" may be used to describe wind power generation or puppy farm.
Farm buildings
Farms have special buildings. Some buildings, such as barns, may hold animals. There may be separate buildings for chickens and pigs. On dairy farms, a milking parlor is an important building. It is where dairy cows are milked. The milk is kept in a milking parlor until a milk tanker comes to get it. There are also special buildings for keeping grain. A silo is a tall building where grains, such as wheat and oats are stored. Farmers also use small round metal buildings to store their grain. These buildings are called grain bins.
Specialized farms
Dairy farm
Dairy farming is a class of agriculture, where female cattle, goats, or other mammals are raised for their milk, which may be either processed on-site or transported to a dairy for processing and eventual retail sale There are many breeds of cattle that can be milked some of the best producing ones include Holstein, Norwegian Red, Kostroma, Brown Swiss, and more.
In most Western countries, a centralized dairy facility processes milk and dairy products, such as ice cream, yogurt, butter, and cheese. In the United States, these dairies are usually local companies, while in the southern hemisphere facilities may be run by very large nationwide or trans-national corporations (such as Fonterra).
Dairy farms generally sell male calves for veal meat, as dairy breeds are not normally satisfactory for commercial beef production. Many dairy farms also grow their own feed, typically including maize, alfalfa, and hay. This is fed directly to the cows, or stored as silage for use during the winter season. Additional dietary supplements are added to the feed to improve milk production.
Poultry farm
Poultry farms are devoted to raising chickens (egg layers or broilers), turkeys, ducks, and other fowl, generally for meat or eggs.
Pig farm
A pig farm is one that specializes in raising pigs for bacon, ham, and other pork products. They may be free range, intensive, or both.
Ownership
Farm control and ownership have traditionally been a key indicator of status and power, especially in Medieval European agrarian societies. The distribution of farm ownership has historically been closely linked to a form of government. Medieval feudalism was essentially a system that centralized control of farmland, control of farm labor, and political power, while the early American democracy, in which land ownership was a prerequisite for voting rights, was built on relatively easy paths to individual farm ownership. However, the gradual modernization and mechanization of farming, which greatly increases both the efficiency and capital requirements of farming, has led to increasingly large farms. This has usually been accompanied by the decoupling of political power from farm ownership.
Forms of ownership
In some societies (especially socialist and communist), collective farming is the norm, with either government ownership of the land or common ownership by a local group. Especially in societies without widespread industrialized farming, tenant farming and sharecropping are common; farmers either pay landowners for the right to use farmland or give up a portion of the crops.
Agribusiness
Farms around the world
United States
The land and buildings of a farm are called the "farmstead".
Enterprises where livestock are raised on rangeland are called ranches. Where livestock are raised in confinement on feed produced elsewhere, the term feedlot is usually used.
In the US, in 1910 there were 6,406,000 farms and 10,174,000 family workers; In 2000 there were only 2,172,000 farms and 2,062,300 family workers. The share of U.S. farms operated by women has risen steadily over recent decades, from 5 percent in 1978 to 14 percent by 2007.In the United States, there are over three million migrant and seasonal farmworkers; 72% are foreign-born, 78% are male, they have an average age of 36 and average education of 8 years. Farmworkers make an average hourly rate of $9–10 per hour, compared to an average of over $18 per hour for nonfarm labor. Their average family income is under $20,000 and 23% live in families with incomes below the federal poverty level. One-half of all farmworker families earn less than $10,000 per year, which is significantly below the 2005 U.S. poverty level of $19,874 for a family of four.
In 2007, corn acres are expected to increase by 15% because of the high demand for ethanol, both in and outside of the U.S. Producers are expecting to plant 90.5 million acres (366,000 km2) of corn, making it the largest corn crop since 1944.
Europe
In the UK, farm as an agricultural unit, always denotes the area of pasture and other fields together with its farmhouse, farmyard and outbuildings. Large farms, or groups of farms under the same ownership, may be called an estate. Conversely, a small farm surrounding the owner's dwelling is called a smallholding and is generally focused on self-sufficiency with only the surplus being sold.
In Europe, traditional family farms are giving way to larger production units where industrial agriculture and mechanization brings large crop yields.
The Common Agricultural Policy (CAP) is one of the most important policies of the European Union and is helping in the change of farms from traditional family farms to larger production units. The policy has the objectives of increasing agricultural production, providing certainty in food supplies, ensuring a high quality of life for farmers, stabilizing markets, and ensuring reasonable prices for consumers. It was, until recently, operated by a system of subsidies and market intervention. Until the 1990s, the policy accounted for over 60 per cent of the European Union's annual budget, and as of 2013 accounts for around 34 per cent.
Asia
Pakistan
According to the World Bank, "most empirical evidence indicates that land productivity on large farms in Pakistan is lower than that of small farms, holding other factors constant." Small farmers have "higher net returns per hectare" than large farms, according to farm
household income data.
Nepal
Nepal is an agricultural country and about 80% of the total population are engaged in farming. Rice is mainly produced in Nepal along with fruits like apples. Dairy farming and poultry farming are also growing in Nepal.
Australia
Farming is a significant economic sector in Australia. A farm is an area of land used for primary production which will include buildings.
According to the UN, "green agriculture directs a greater share of total farming input expenditures towards the purchase of locally sourced input?(e.g. labour and organic fertilisers) and a local multiplier effect is expected to kick in. Overall, green farming practices tend to require more labour inputs than conventional farming (e.g. from comparable levels to as much as 30 percent more) (FAO 2007 and European Commission 2010), creating jobs in rural areas and a higher return on labour inputs."
Often very small farms used for intensive primary production are referred to by the specialization they are being used for, such as a dairy rather than a dairy farm, a piggery, a market garden, etc. This also applies to feedlots, which are specifically developed for a single purpose and are often not able to be used for more general purpose (mixed) farming practices.
In remote areas, farms can become quite large. As with estates in England, there is no defined size or method of operation at which a large farm becomes a station.
Africa
A farm in Africa includes various structures. Depending on climate-related areas primarily farming is the raising and breeding of grazing livestock, such as cattle, sheep, ostriches, horses or goats. Predominantly domestic animals are raised for their meat, milk, skin, leather or fiber wool). You might even come across silk farms.
Furthermore, there are plenty of hunting farms, guest farms and game farms. Arable] or irrigated land is often used for raising crops such as feed grains and hay for animal feeding.
On some farms (Astro Farm) star-gazing became very popular because of the excellent optical quality in the desert. The High Energy Stereoscopic System (H.E.S.S.) which investigates cosmic gamma rays is situated on Farm Göllschau in Namibia.
Farm equipment
Farm equipment has evolved over the centuries from simple hand tools such as the hoe, through ox- or horse-drawn equipment such as the plough and harrow, to the modern highly technical machinery such as the tractor, baler and combine harvester replacing what was a highly labour-intensive occupation before the Industrial Revolution. Today much of the farm equipment used on both small and large farms is automated (e.g. using satellite guided farming).
As new types of high-tech farm equipment have become inaccessible to farmers that historically fixed their own equipment, Wired magazine reports there is a growing backlash, due mostly to companies using intellectual property law to prevent farmers from having the legal right to fix their equipment (or gain access to the information to allow them to do it). This has encouraged groups such as Open Source Ecology and Farm Hack to begin to make open source hardware for agricultural machinery. In addition on a smaller scale Farmbot and the RepRap open source 3D printer community has begun to make open-source farm tools available of increasing levels of sophistication.
| Technology | Buildings and infrastructure | null |
59856 | https://en.wikipedia.org/wiki/Rhea%20%28moon%29 | Rhea (moon) | Rhea () is the second-largest moon of Saturn and the ninth-largest moon in the Solar System, with a surface area that is comparable to the area of Australia. It is the smallest body in the Solar System for which precise measurements have confirmed a shape consistent with hydrostatic equilibrium. Rhea has a nearly circular orbit around Saturn, but it is also tidally locked, like Saturn's other major moons; that is, it rotates with the same period it revolves (orbits), so one hemisphere always faces towards the planet.
The moon itself has a fairly low density, composed of roughly three-quarters ice and only one-quarter rock. The surface of Rhea is heavily cratered, with distinct leading and trailing hemispheres. Like the moon Dione, it has high-albedo ice cliffs that appear as bright wispy streaks visible from space. The surface temperature varies between −174 °C and −220 °C.
Rhea was discovered in 1672 by Giovanni Domenico Cassini. Since then, it has been visited by both Voyager probes and was the subject of close targeted flybys by the Cassini orbiter in 2005, 2007, 2010, 2011, and once more in 2013.
Discovery
Rhea was discovered by Giovanni Domenico Cassini on 23 December 1672, with a telescope made by Giuseppe Campani. Cassini named the four moons he discovered (Tethys, Dione, Rhea, and Iapetus) Sidera Lodoicea (the stars of Louis) to honor King Louis XIV. Rhea was the second moon of Saturn that Cassini discovered, and the third moon discovered around Saturn overall.
Name
Rhea is named after the Titan Rhea of Greek mythology, the mother of the first generation of Olympian gods and wife of Cronus, the Greek counterpart of the god Saturn. It is also designated Saturn V (being the fifth major moon going outward from the planet, after Mimas, Enceladus, Tethys, and Dione).
Astronomers fell into the habit of referring to them and Titan as Saturn I through Saturn V. Once Mimas and Enceladus were discovered, in 1789, the numbering scheme was extended to Saturn VII, and then to Saturn VIII with the discovery of Hyperion in 1848.
Rhea was not named until 1847, when John Herschel (son of William Herschel, discoverer of the planet Uranus and two other moons of Saturn, Mimas and Enceladus) suggested in Results of Astronomical Observations made at the Cape of Good Hope that the names of the Titans, sisters and brothers of Kronos (Saturn, in Roman mythology), be used.
Orbit
The orbit of Rhea has very low eccentricity (0.001), meaning it is nearly circular. It has a low inclination of less than a degree, inclined by only 0.35° from Saturn's equatorial plane.
Rhea is tidally locked and rotates synchronously; that is, it rotates at the same speed it revolves (orbits), so one hemisphere is always facing towards Saturn. This is called the near pole. Equally, one hemisphere always faces forward, relative to the direction of movement; this is called the leading hemisphere; the other side is the trailing hemisphere, which faces backwards relative to the moon's motion.
Physical characteristics
Size, mass, and internal structure
Rhea is the second largest moon of Saturn, but with a mean diameter of 1,528 kilometers (949 miles) it is less than a third the radius of Saturn's largest moon, Titan. Rhea is an icy body with a density of about 1.236 g/cm3. This low density indicates that it is made of ~25% rock (density ~3.25 g/cm3) and ~75% water ice (density ~0.93 g/cm3). A layer of Ice II (a high-pressure and extra-low temperature form of ice) is believed, based on the moon's temperature profile, to start around beneath the surface. Saturn's biggest moon. Although Rhea is the ninth-largest moon, it is only the tenth-most massive moon. Indeed, Oberon, the second-largest moon of Uranus, has almost the same size, but is significantly denser than Rhea (1.63 vs 1.24) and thus more massive, although Rhea is slightly larger by volume. The surface area of the moon can be estimated at , similar to Australia (7,688,287 km2).
Before the Cassini–Huygens mission, it was assumed that Rhea had a rocky core. However, measurements taken during a close flyby by the Cassini orbiter in 2005 cast this into doubt. In a paper published in 2007 it was claimed that the axial dimensionless moment of inertia coefficient was 0.4. Such a value indicated that Rhea had an almost homogeneous interior (with some compression of ice in the center) while the existence of a rocky core would imply a moment of inertia of about 0.34. In the same year, another paper claimed the moment of inertia was about 0.37. Rhea being either partially or fully differentiated would be consistent with the observations of the Cassini probe. A year later, yet another paper claimed that the moon may not be in hydrostatic equilibrium, meaning that the moment of inertia cannot be determined from the gravity data alone. In 2008, an author of the first paper tried to reconcile these three disparate results. He concluded that there is a systematic error in the Cassini radio Doppler data used in the analysis, but, after restricting the analysis to a subset of data obtained closest to the moon, he arrived at his old result that Rhea was in hydrostatic equilibrium and had a moment of inertia of about 0.4, again implying a homogeneous interior.
The triaxial shape of Rhea is consistent with a homogeneous body in hydrostatic equilibrium rotating at Rhea's angular velocity. Modelling in 2006 suggested that Rhea could be barely capable of sustaining an internal liquid-water ocean through heating by radioactive decay; such an ocean would have to be at about 176 K, the eutectic temperature for the water–ammonia system. More recent indications are that Rhea has a homogeneous interior and hence that this ocean does not exist.
Surface features
Rhea's features resemble those of Dione, with distinct and dissmillar leading and trailing hemispheres, suggesting similar composition and histories. The temperature on Rhea is 99 K (−174 °C) in direct sunlight and between 73 K (−200 °C) and 53 K (−220 °C) in the shade.
Rhea has a rather typical heavily cratered surface, with the exceptions of a few large Dione-type chasmata or fractures (formerly known as wispy terrain) on the trailing hemisphere (the side facing away from the direction of motion along Rhea's orbit) and a very faint "line" of material at Rhea's equator that may have been deposited by material deorbiting from its rings. Rhea has two very large impact basins on its hemisphere facing away from Saturn, which are about 400 and 500 km across. The more northerly and less degraded of the two, called Tirawa, is roughly comparable in size to the basin Odysseus on Tethys. There is a 48 km-diameter impact crater at 112°W that is prominent because of an extended system of bright rays, which extend up to away from the crater, across most of one hemisphere. This crater, called Inktomi, is nicknamed "The Splat", and may be one of the youngest craters on the inner moons of Saturn. This was hypothesized in a 2007 paper published by Lunar and Planetary Science. Rhea's impact craters are more crisply defined than the flatter craters that are pervasive on Ganymede and Callisto; it is theorized that this is due to a much lower surface gravity (0.26 m/s2, compared to Ganymede's 1.428 m/s2 and Callisto's 1.235 m/s2) and a stiffer crust of ice. Similarly, ejecta blankets – asymmetrical blankets of ejected particles surrounding impact craters – are not present on Rhea, potentially another result of the moon's low surface gravity.
Its surface can be divided into two geologically different areas based on crater density; the first area contains craters which are larger than 40 km in diameter, whereas the second area, in parts of the polar and equatorial regions, has only craters under that size. This suggests that a major resurfacing event occurred some time during its formation. The leading hemisphere is heavily cratered and uniformly bright. As on Callisto, the craters lack the high relief features seen on the Moon and Mercury. It has been theorized that these cratered plains are up to four billion years old on average. On the trailing hemisphere there is a network of bright swaths on a dark background, and fewer craters. It is believed, based on data from the Cassini probe, that these are tectonic features: depressions (graben) and troughs, with ice-covered cliff sides causing the lines' whiteness (more technically their albedo). The extensive dark areas are thought to be deposited tholins, which are a mix of complex organic compounds generated on the ice by pyrolysis and radiolysis of simple compounds containing carbon, nitrogen and hydrogen. The trailing side of Rhea's surface is irradiated by Saturn's magnetosphere, which may cause chemical-level changes on the surface, including radiolysis (see ). Particles from Saturn's E-ring are also flung onto the moon's leading hemisphere, coating it.
Rhea has some evidence of endogenic activity – that is, activity originating from within the moon, such as heating and cryovolcanic activity: there are fault systems and craters with uplifted bases (so-called "relaxed" craters), although the latter is apparently only present in large craters more than across.
Formation
The moons of Saturn are thought to have formed through co-accretion, a similar process to that believed to have formed the planets in the Solar System. As the young giant planets formed, they were surrounded by discs of material that gradually coalesced into moons. However, a model proposed by Erik Asphaug and Andreas Reufer for the formation of Titan may also shine a new light on the origin of Rhea and Iapetus. In this model, Titan was formed in a series of giant impacts between pre-existing moons, and Rhea and Iapetus are thought to have formed from part of the debris of these collisions.
Atmosphere
On November 27, 2010, NASA announced the discovery of an extremely tenuous atmosphere—an exosphere. It consists of oxygen and carbon dioxide in proportion of roughly 5 to 2. The surface density of the exosphere is from 105 to 106 molecules in a cubic centimeter, depending on local temperature. The main source of oxygen is radiolysis of water ice at the surface via irradiation from the magnetosphere of Saturn. The source of the carbon dioxide is less clear, but it may be related to oxidation of the organics present in ice or to outgassing of the moon's interior.
Possible ring system
On March 6, 2008, NASA announced that Rhea may have a weak ring system. This would mark the first discovery of rings around a moon. The rings' existence was inferred by observed changes in the flow of electrons trapped by Saturn's magnetic field as Cassini passed by Rhea. Dust and debris could extend out to Rhea's Hill sphere, but were thought to be denser nearer the moon, with three narrow rings of higher density. The case for a ring was strengthened by the subsequent finding of the presence of a set of small ultraviolet-bright spots distributed along Rhea's equator (interpreted as the impact points of deorbiting ring material). However, when Cassini made targeted observations of the putative ring plane from several angles, there was no evidence of ring material found, suggesting that another explanation for the earlier observations is needed.
Exploration
The first images of Rhea were obtained by Voyager 1 & 2 spacecraft in 1980–1981.
There were five close targeted fly-bys by the Cassini orbiter, which was one part of the dual orbiter and lander Cassini–Huygens mission. Launched in 1997, Cassini–Huygens was targeted at the Saturn system; in total it took more than 450 thousand images. Cassini passed Rhea at a distance of 500 km on November 26, 2005; at a distance of 5,750 km on August 30, 2007; at a distance of 100 km on March 2, 2010; at 69 km flyby on January 11, 2011; and a last flyby at 992 km on March 9, 2013.
| Physical sciences | Solar System | Astronomy |
59858 | https://en.wikipedia.org/wiki/Rhea%20%28bird%29 | Rhea (bird) | Rheas ( ), also known as ñandus ( ) or South American ostrich, are moderately sized South American ratites (flightless birds without a keel on their sternum bone) of the order Rheiformes. They are distantly related to the African ostriches and Australia's emu (the largest and second-largest living ratites, respectively), with rheas placing just behind the emu in height and overall size.
Most taxonomic authorities recognize two extant species: the greater or American rhea (Rhea americana), and the lesser or Darwin's rhea (Rhea pennata). The International Union for Conservation of Nature (IUCN) classifies the puna rhea as another species instead of a subspecies of the lesser rhea. The IUCN currently rates the greater and puna rheas as near-threatened in their native ranges, while Darwin's rhea is of least concern, having recovered from past threats to its survival. In addition, the feral population of the greater rhea in Germany appears to be growing. However, control efforts are underway and seem to succeed in controlling the birds' population growth. Similarly to ostriches and emus, rheas are fairly popular livestock and pets, regularly kept and bred on farms, ranches, private parks, and by aviculturists, mainly in North and South America and Europe.
Etymology
The name "rhea" was used in 1752 by Paul Möhring and adopted as the English common name. Möhring named the rhea after the Greek Titan Rhea, whose Ancient Greek name () is thought to come from (éra, "ground"). This was fitting with the rhea being a flightless ground bird. Depending on the South American region, the rhea is known locally as (Guaraní –or related Tupi nhandú-gûasú– meaning "big spider" most probably concerning their habit of opening and lowering alternate wings when they run), (Portuguese), (Aymara and Quechua), or (Mapudungun). is the common name in many European languages and may sometimes be heard in English.
Taxonomy and systematics
The genus Rhea was introduced by French zoologist Mathurin Jacques Brisson in 1760 with the greater rhea (Rhea americana) as the type species.
Extant species
The genus contains two extant species and eight subspecies, although one subspecies is disputed:
Rhea pennata was not always in the genus Rhea. In 2008, the SACC, the last holdout, approved merging the genera Rhea and Pterocnemia on August 7, 2008. This merging of genera leaves only the genus Rhea. A former fourth species of rhea, Rhea nana, was described by Lydekker in 1894 based on a single egg found in Patagonia, but today no major authorities consider it valid.
Fossils
†R. anchorenense (Ameghino & Rusconi 1932) [Rhea americana anchorenense Amcghino & Rusconi 1932]
†R. fossilis (Moreno & Mercerat 1891) [Pterocnemia fossilis (Moreno & Mercerat 1891); Rhea pampeana (Moreno & Mercerat 1891)]
†R. mesopotamica (Agnolín & Noriega 2012) [Pterocnemia mesopotamica Agnolín & Noriega 2012]
†R. subpampeana Moreno & Mercerat 1891
Description
Rheas are large, flightless birds with grey-brown plumage, long legs, and long necks, similar to an ostrich. Large males of R. americana can reach tall at the head, at the back and can weigh up to . The lesser rhea is smaller, with a height of . Their wings are large for a flightless bird () and are spread while running, to act like sails. Unlike most birds, rheas have only three toes. Their tarsus has 18 to 22 horizontal plates on the front of it. They also store urine separately in an expansion of the cloaca.
Distribution and habitat
Rheas are from South America only and are limited within the continent to Argentina, Bolivia, Brazil, Chile, Paraguay, Peru and Uruguay. They are grassland birds, and both species prefer open land. The greater rheas live in open grasslands, pampas and chaco woodlands. They prefer to breed near water and prefer lowlands, seldom going above . On the other hand, the lesser rhea will inhabit most shrubland, grassland, even desert salt puna up to .
Feral populations in Europe
A small population of rheas has emerged in Mecklenburg-Western Pomerania, northeastern Germany, after several couples escaped from an exotic meat farm near Lübeck in the late 1990s. Contrary to expectations, the large birds adapted well to conditions in the German countryside. A monitoring system has been in place since 2008. By 2014, there was already a population of well over 100 birds in an area of between the river Wakenitz and the A20 motorway, slowly expanding eastward.
The population grew steadily for several years. By autumn 2018, their numbers had significantly increased to about 600. As such, local farmers claim increasing damage to their fields, and some biologists say the rheas pose a growing risk to local wildlife. Still protected by German natural conservation law, a local discussion developed regarding how to handle the situation. Eventually, Mecklenburg-Western Pomerania's government allowed limited hunting of the birds, explicitly to just reduce the population's growth and not to wipe them out. At this point, it was generally agreed that the rheas should be allowed to stay in the region. By spring 2021, just 247 rheas were counted; this development was attributed to both the hunting and the increased caution of the animals. Several had begun to avoid humans more than previously and retreated into the woods. Some members of this rhea population have also expanded into other areas; at least twice individual rheas who probably originated in Mecklenburg-Western Pomerania were sighted in Brandenburg's High Fläming Nature Park, over from their usual range. By early 2023, 91 rheas were counted in Mecklenburg-Western Pomerania; the population decline was attributed to both hunting as well as harsher weather of previous years. By this point, German authorities believed a stable population of 50 adult birds would be optimal for the local ecosystem and agriculture. Researchers concluded that the feral population was subject to substantial fluctuations but remained healthy, adaptable, and entrenched in the area.
There also appears to be a small population of wild rheas in the United Kingdom. In March 2021, about 20 rheas were reportedly running free on a residential estate in Hertfordshire. Local police could not identify any owner, so they assumed they were wild birds. Once caught, authorities intend to place them in a suitable nature reserve to allow them to develop as a colony.
Behavior
Individual and flocking
Rheas tend to be silent birds, except when they are chicks or the male seeks a mate. During the breeding season, the male will attempt to attract females by calling. This call is a loud booming noise. While calling like this, they will lift the front of their body and ruffle their plumage, all while keeping their neck stiff. They will then extend and raise their wings and run short distances, alternating with their wings. He may then single out a female and walk alongside or in front of her with a lowered head and spread wings. If the female notices him, he will wave his neck back and forth in a figure eight. Finally, a female may offer herself, and copulation will commence.
During the non-breeding season they may form flocks of between 20 and 25 birds, although the lesser rhea forms smaller flocks than this. When in danger, they flee in a zigzag course, using one wing and the other, similar to a rudder. During the breeding season, the flocks break up.
Diet
Mostly, rheas are herbivorous and prefer broad-leafed plants, but they also eat fruits, seeds, roots, and insects such as grasshoppers, small reptiles, and rodents. Young rheas eat only insects for the first few days. Outside the breeding season, they gather in flocks and feed with deer and cattle.
Reproduction
Rheas are polygynandrous, with males courting between two and twelve females and females commonly mating with multiple dominant males during the breeding season. After mating, the male builds a nest where each female lays eggs. The nest is a simple scrape in the ground, lined with grass and leaves. The male incubates from ten to sixty eggs. The male will use a decoy system and place some eggs outside the nest, then sacrifice these to predators so they do not attempt to get inside the nest. The male may use another subordinate male to incubate his eggs while he finds another group of females to start a second nest with. The chicks hatch within 36 hours of each other. Right before hatching, the chicks begin to whistle. The group of females, meanwhile, may move on and mate with other males. While caring for the young, the males will charge at any perceived threat approaching the chicks, including female rheas and humans. The young reach full adult size in about six months but do not breed until they reach two years of age.
Status and conservation
The numbers of the greater and puna rhea are decreasing as their habitats shrink. Both are considered near threatened by the IUCN. The IUCN also states that they are both approaching vulnerable status. The lesser rhea is classified as least concern.
Human interaction
Rheas have many uses in South America. Feathers are used for feather dusters, skins are used for cloaks or leather, and their meat is a staple to many people.
Gauchos traditionally hunt rheas on horseback, throwing bolas or boleadoras—a throwing device consisting of three balls joined by rope—at their legs, which immobilises the bird. The rhea is pictured on Argentina's 1-centavo coin minted in 1987, and on the Uruguayan 5-peso coin.
| Biology and health sciences | Ratites | null |
59861 | https://en.wikipedia.org/wiki/Experiment | Experiment | An experiment is a procedure carried out to support or refute a hypothesis, or determine the efficacy or likelihood of something previously untried. Experiments provide insight into cause-and-effect by demonstrating what outcome occurs when a particular factor is manipulated. Experiments vary greatly in goal and scale but always rely on repeatable procedure and logical analysis of the results. There also exist natural experimental studies.
A child may carry out basic experiments to understand how things fall to the ground, while teams of scientists may take years of systematic investigation to advance their understanding of a phenomenon. Experiments and other types of hands-on activities are very important to student learning in the science classroom. Experiments can raise test scores and help a student become more engaged and interested in the material they are learning, especially when used over time. Experiments can vary from personal and informal natural comparisons (e.g. tasting a range of chocolates to find a favorite), to highly controlled (e.g. tests requiring complex apparatus overseen by many scientists that hope to discover information about subatomic particles). Uses of experiments vary considerably between the natural and human sciences.
Experiments typically include controls, which are designed to minimize the effects of variables other than the single independent variable. This increases the reliability of the results, often through a comparison between control measurements and the other measurements. Scientific controls are a part of the scientific method. Ideally, all variables in an experiment are controlled (accounted for by the control measurements) and none are uncontrolled. In such an experiment, if all controls work as expected, it is possible to conclude that the experiment works as intended, and that results are due to the effect of the tested variables.
Overview
In the scientific method, an experiment is an empirical procedure that arbitrates competing models or hypotheses. Researchers also use experimentation to test existing theories or new hypotheses to support or disprove them.
An experiment usually tests a hypothesis, which is an expectation about how a particular process or phenomenon works. However, an experiment may also aim to answer a "what-if" question, without a specific expectation about what the experiment reveals, or to confirm prior results. If an experiment is carefully conducted, the results usually either support or disprove the hypothesis. According to some philosophies of science, an experiment can never "prove" a hypothesis, it can only add support. On the other hand, an experiment that provides a counterexample can disprove a theory or hypothesis, but a theory can always be salvaged by appropriate ad hoc modifications at the expense of simplicity.
An experiment must also control the possible confounding factors—any factors that would mar the accuracy or repeatability of the experiment or the ability to interpret the results. Confounding is commonly eliminated through scientific controls and/or, in randomized experiments, through random assignment.
In engineering and the physical sciences, experiments are a primary component of the scientific method. They are used to test theories and hypotheses about how physical processes work under particular conditions (e.g., whether a particular engineering process can produce a desired chemical compound). Typically, experiments in these fields focus on replication of identical procedures in hopes of producing identical results in each replication. Random assignment is uncommon.
In medicine and the social sciences, the prevalence of experimental research varies widely across disciplines. When used, however, experiments typically follow the form of the clinical trial, where experimental units (usually individual human beings) are randomly assigned to a treatment or control condition where one or more outcomes are assessed. In contrast to norms in the physical sciences, the focus is typically on the average treatment effect (the difference in outcomes between the treatment and control groups) or another test statistic produced by the experiment. A single study typically does not involve replications of the experiment, but separate studies may be aggregated through systematic review and meta-analysis.
There are various differences in experimental practice in each of the branches of science. For example, agricultural research frequently uses randomized experiments (e.g., to test the comparative effectiveness of different fertilizers), while experimental economics often involves experimental tests of theorized human behaviors without relying on random assignment of individuals to treatment and control conditions.
History
One of the first methodical approaches to experiments in the modern sense is visible in the works of the Arab mathematician and scholar Ibn al-Haytham. He conducted his experiments in the field of optics—going back to optical and mathematical problems in the works of Ptolemy—by controlling his experiments due to factors such as self-criticality, reliance on visible results of the experiments as well as a criticality in terms of earlier results. He was one of the first scholars to use an inductive-experimental method for achieving results. In his Book of Optics he describes the fundamentally new approach to knowledge and research in an experimental sense:
According to his explanation, a strictly controlled test execution with a sensibility for the subjectivity and susceptibility of outcomes due to the nature of man is necessary. Furthermore, a critical view on the results and outcomes of earlier scholars is necessary:
Thus, a comparison of earlier results with the experimental results is necessary for an objective experiment—the visible results being more important. In the end, this may mean that an experimental researcher must find enough courage to discard traditional opinions or results, especially if these results are not experimental but results from a logical/ mental derivation. In this process of critical consideration, the man himself should not forget that he tends to subjective opinions—through "prejudices" and "leniency"—and thus has to be critical about his own way of building hypotheses.
Francis Bacon (1561–1626), an English philosopher and scientist active in the 17th century, became an influential supporter of experimental science in the English renaissance. He disagreed with the method of answering scientific questions by deduction—similar to Ibn al-Haytham—and described it as follows: "Having first determined the question according to his will, man then resorts to experience, and bending her to conformity with his placets, leads her about like a captive in a procession." Bacon wanted a method that relied on repeatable observations, or experiments. Notably, he first ordered the scientific method as we understand it today.
In the centuries that followed, people who applied the scientific method in different areas made important advances and discoveries. For example, Galileo Galilei (1564–1642) accurately measured time and experimented to make accurate measurements and conclusions about the speed of a falling body. Antoine Lavoisier (1743–1794), a French chemist, used experiment to describe new areas, such as combustion and biochemistry and to develop the theory of conservation of mass (matter). Louis Pasteur (1822–1895) used the scientific method to disprove the prevailing theory of spontaneous generation and to develop the germ theory of disease. Because of the importance of controlling potentially confounding variables, the use of well-designed laboratory experiments is preferred when possible.
A considerable amount of progress on the design and analysis of experiments occurred in the early 20th century, with contributions from statisticians such as Ronald Fisher (1890–1962), Jerzy Neyman (1894–1981), Oscar Kempthorne (1919–2000), Gertrude Mary Cox (1900–1978), and William Gemmell Cochran (1909–1980), among others.
Types
Experiments might be categorized according to a number of dimensions, depending upon professional norms and standards in different fields of study.
In some disciplines (e.g., psychology or political science), a 'true experiment' is a method of social research in which there are two kinds of variables. The independent variable is manipulated by the experimenter, and the dependent variable is measured. The signifying characteristic of a true experiment is that it randomly allocates the subjects to neutralize experimenter bias, and ensures, over a large number of iterations of the experiment, that it controls for all confounding factors.
Depending on the discipline, experiments can be conducted to accomplish different but not mutually exclusive goals: test theories, search for and document phenomena, develop theories, or advise policymakers. These goals also relate differently to validity concerns.
Controlled experiments
A controlled experiment often compares the results obtained from experimental samples against control samples, which are practically identical to the experimental sample except for the one aspect whose effect is being tested (the independent variable). A good example would be a drug trial. The sample or group receiving the drug would be the experimental group (treatment group); and the one receiving the placebo or regular treatment would be the control one. In many laboratory experiments it is good practice to have several replicate samples for the test being performed and have both a positive control and a negative control. The results from replicate samples can often be averaged, or if one of the replicates is obviously inconsistent with the results from the other samples, it can be discarded as being the result of an experimental error (some step of the test procedure may have been mistakenly omitted for that sample). Most often, tests are done in duplicate or triplicate. A positive control is a procedure similar to the actual experimental test but is known from previous experience to give a positive result. A negative control is known to give a negative result. The positive control confirms that the basic conditions of the experiment were able to produce a positive result, even if none of the actual experimental samples produce a positive result. The negative control demonstrates the base-line result obtained when a test does not produce a measurable positive result. Most often the value of the negative control is treated as a "background" value to subtract from the test sample results. Sometimes the positive control takes the quadrant of a standard curve.
An example that is often used in teaching laboratories is a controlled protein assay. Students might be given a fluid sample containing an unknown (to the student) amount of protein. It is their job to correctly perform a controlled experiment in which they determine the concentration of protein in the fluid sample (usually called the "unknown sample"). The teaching lab would be equipped with a protein standard solution with a known protein concentration. Students could make several positive control samples containing various dilutions of the protein standard. Negative control samples would contain all of the reagents for the protein assay but no protein. In this example, all samples are performed in duplicate. The assay is a colorimetric assay in which a spectrophotometer can measure the amount of protein in samples by detecting a colored complex formed by the interaction of protein molecules and molecules of an added dye. In the illustration, the results for the diluted test samples can be compared to the results of the standard curve (the blue line in the illustration) to estimate the amount of protein in the unknown sample.
Controlled experiments can be performed when it is difficult to exactly control all the conditions in an experiment. In this case, the experiment begins by creating two or more sample groups that are probabilistically equivalent, which means that measurements of traits should be similar among the groups and that the groups should respond in the same manner if given the same treatment. This equivalency is determined by statistical methods that take into account the amount of variation between individuals and the number of individuals in each group. In fields such as microbiology and chemistry, where there is very little variation between individuals and the group size is easily in the millions, these statistical methods are often bypassed and simply splitting a solution into equal parts is assumed to produce identical sample groups.
Once equivalent groups have been formed, the experimenter tries to treat them identically except for the one variable that he or she wishes to isolate. Human experimentation requires special safeguards against outside variables such as the placebo effect. Such experiments are generally double blind, meaning that neither the volunteer nor the researcher knows which individuals are in the control group or the experimental group until after all of the data have been collected. This ensures that any effects on the volunteer are due to the treatment itself and are not a response to the knowledge that he is being treated.
In human experiments, researchers may give a subject (person) a stimulus that the subject responds to. The goal of the experiment is to measure the response to the stimulus by a test method.
In the design of experiments, two or more "treatments" are applied to estimate the difference between the mean responses for the treatments. For example, an experiment on baking bread could estimate the difference in the responses associated with quantitative variables, such as the ratio of water to flour, and with qualitative variables, such as strains of yeast. Experimentation is the step in the scientific method that helps people decide between two or more competing explanations—or hypotheses. These hypotheses suggest reasons to explain a phenomenon or predict the results of an action. An example might be the hypothesis that "if I release this ball, it will fall to the floor": this suggestion can then be tested by carrying out the experiment of letting go of the ball, and observing the results. Formally, a hypothesis is compared against its opposite or null hypothesis ("if I release this ball, it will not fall to the floor"). The null hypothesis is that there is no explanation or predictive power of the phenomenon through the reasoning that is being investigated. Once hypotheses are defined, an experiment can be carried out and the results analysed to confirm, refute, or define the accuracy of the hypotheses.
Experiments can be also designed to estimate spillover effects onto nearby untreated units.
Natural experiments
The term "experiment" usually implies a controlled experiment, but sometimes controlled experiments are prohibitively difficult, impossible, unethical or illegal. In this case researchers resort to natural experiments or quasi-experiments. Natural experiments rely solely on observations of the variables of the system under study, rather than manipulation of just one or a few variables as occurs in controlled experiments. To the degree possible, they attempt to collect data for the system in such a way that contribution from all variables can be determined, and where the effects of variation in certain variables remain approximately constant so that the effects of other variables can be discerned. The degree to which this is possible depends on the observed correlation between explanatory variables in the observed data. When these variables are not well correlated, natural experiments can approach the power of controlled experiments. Usually, however, there is some correlation between these variables, which reduces the reliability of natural experiments relative to what could be concluded if a controlled experiment were performed. Also, because natural experiments usually take place in uncontrolled environments, variables from undetected sources are neither measured nor held constant, and these may produce illusory correlations in variables under study.
Much research in several science disciplines, including economics, human geography, archaeology, sociology, cultural anthropology, geology, paleontology, ecology, meteorology, and astronomy, relies on quasi-experiments. For example, in astronomy it is clearly impossible, when testing the hypothesis "Stars are collapsed clouds of hydrogen", to start out with a giant cloud of hydrogen, and then perform the experiment of waiting a few billion years for it to form a star. However, by observing various clouds of hydrogen in various states of collapse, and other implications of the hypothesis (for example, the presence of various spectral emissions from the light of stars), we can collect data we require to support the hypothesis. An early example of this type of experiment was the first verification in the 17th century that light does not travel from place to place instantaneously, but instead has a measurable speed. Observation of the appearance of the moons of Jupiter were slightly delayed when Jupiter was farther from Earth, as opposed to when Jupiter was closer to Earth; and this phenomenon was used to demonstrate that the difference in the time of appearance of the moons was consistent with a measurable speed.
Field experiments
Field experiments are so named to distinguish them from laboratory experiments, which enforce scientific control by testing a hypothesis in the artificial and highly controlled setting of a laboratory. Often used in the social sciences, and especially in economic analyses of education and health interventions, field experiments have the advantage that outcomes are observed in a natural setting rather than in a contrived laboratory environment. For this reason, field experiments are sometimes seen as having higher external validity than laboratory experiments. However, like natural experiments, field experiments suffer from the possibility of contamination: experimental conditions can be controlled with more precision and certainty in the lab. Yet some phenomena (e.g., voter turnout in an election) cannot be easily studied in a laboratory.
Observational studies
An observational study is used when it is impractical, unethical, cost-prohibitive (or otherwise inefficient) to fit a physical or social system into a laboratory setting, to completely control confounding factors, or to apply random assignment. It can also be used when confounding factors are either limited or known well enough to analyze the data in light of them (though this may be rare when social phenomena are under examination). For an observational science to be valid, the experimenter must know and account for confounding factors. In these situations, observational studies have value because they often suggest hypotheses that can be tested with randomized experiments or by collecting fresh data.
Fundamentally, however, observational studies are not experiments. By definition, observational studies lack the manipulation required for Baconian experiments. In addition, observational studies (e.g., in biological or social systems) often involve variables that are difficult to quantify or control. Observational studies are limited because they lack the statistical properties of randomized experiments. In a randomized experiment, the method of randomization specified in the experimental protocol guides the statistical analysis, which is usually specified also by the experimental protocol. Without a statistical model that reflects an objective randomization, the statistical analysis relies on a subjective model. Inferences from subjective models are unreliable in theory and practice. In fact, there are several cases where carefully conducted observational studies consistently give wrong results, that is, where the results of the observational studies are inconsistent and also differ from the results of experiments. For example, epidemiological studies of colon cancer consistently show beneficial correlations with broccoli consumption, while experiments find no benefit.
A particular problem with observational studies involving human subjects is the great difficulty attaining fair comparisons between treatments (or exposures), because such studies are prone to selection bias, and groups receiving different treatments (exposures) may differ greatly according to their covariates (age, height, weight, medications, exercise, nutritional status, ethnicity, family medical history, etc.). In contrast, randomization implies that for each covariate, the mean for each group is expected to be the same. For any randomized trial, some variation from the mean is expected, of course, but the randomization ensures that the experimental groups have mean values that are close, due to the central limit theorem and Markov's inequality. With inadequate randomization or low sample size, the systematic variation in covariates between the treatment groups (or exposure groups) makes it difficult to separate the effect of the treatment (exposure) from the effects of the other covariates, most of which have not been measured. The mathematical models used to analyze such data must consider each differing covariate (if measured), and results are not meaningful if a covariate is neither randomized nor included in the model.
To avoid conditions that render an experiment far less useful, physicians conducting medical trials—say for U.S. Food and Drug Administration approval—quantify and randomize the covariates that can be identified. Researchers attempt to reduce the biases of observational studies with matching methods such as propensity score matching, which require large populations of subjects and extensive information on covariates. However, propensity score matching is no longer recommended as a technique because it can increase, rather than decrease, bias. Outcomes are also quantified when possible (bone density, the amount of some cell or substance in the blood, physical strength or endurance, etc.) and not based on a subject's or a professional observer's opinion. In this way, the design of an observational study can render the results more objective and therefore, more convincing.
Ethics
By placing the distribution of the independent variable(s) under the control of the researcher, an experiment—particularly when it involves human subjects—introduces potential ethical considerations, such as balancing benefit and harm, fairly distributing interventions (e.g., treatments for a disease), and informed consent. For example, in psychology or health care, it is unethical to provide a substandard treatment to patients. Therefore, ethical review boards are supposed to stop clinical trials and other experiments unless a new treatment is believed to offer benefits as good as current best practice. It is also generally unethical (and often illegal) to conduct randomized experiments on the effects of substandard or harmful treatments, such as the effects of ingesting arsenic on human health. To understand the effects of such exposures, scientists sometimes use observational studies to understand the effects of those factors.
Even when experimental research does not directly involve human subjects, it may still present ethical concerns. For example, the nuclear bomb experiments conducted by the Manhattan Project implied the use of nuclear reactions to harm human beings even though the experiments did not directly involve any human subjects.
| Physical sciences | Basics | null |
59863 | https://en.wikipedia.org/wiki/Correspondence%20principle | Correspondence principle | In physics, a correspondence principle is any one of several premises or assertions about the relationship between classical and quantum mechanics.
The physicist Niels Bohr coined the term in 1920 during the early development of quantum theory; he used it to explain how quantized classical orbitals connect to quantum radiation.
Modern sources often use the term for the idea that the behavior of systems described by quantum theory reproduces classical physics in the limit of large quantum numbers: for large orbits and for large energies, quantum calculations must agree with classical calculations. A "generalized" correspondence principle refers to the requirement for a broad set of connections between any old and new theory.
History
Max Planck was the first to introduce the idea of quanta of energy, while studying black-body radiation in 1900. In 1906, he was also the first to write that quantum theory should replicate classical mechanics at some limit, particularly if the Planck constant h were taken to be infinitesimal. With this idea, he showed that Planck's law for thermal radiation leads to the Rayleigh–Jeans law, the classical prediction (valid for large wavelength).
Niels Bohr used a similar idea, while developing his model of the atom. In 1913, he provided the first postulates of what is now known as old quantum theory. Using these postulates he obtained that for the hydrogen atom, the energy spectrum approaches the classical continuum for large n (a quantum number that encodes the energy of the orbit). Bohr coined the term "correspondence principle" during a lecture in 1920.
Arnold Sommerfeld refined Bohr's theory leading to the Bohr-Sommerfeld quantization condition. Sommerfeld referred to the correspondence principle as Bohr's magic wand (), in 1921.
Bohr's correspondence principle
The seeds of Bohr's correspondence principle appeared from two sources. First Sommerfeld and Max Born developed a "quantization procedure" based on the action angle variables of classical Hamiltonian mechanics. This gave a mathematical foundation for stationary states of the Bohr-Sommerfeld model of the atom. The second seed was Albert Einstein's quantum derivation of Planck's law in 1916. Einstein developed the statistical mechanics for Bohr-model atoms interacting with electromagnetic radiation, leading to absorption and two kinds of emission, spontaneous and stimulated emission. But for Bohr the important result was the use of classical analogies and the Bohr atomic model to fix inconsistencies in Planck's derivation of the blackbody radiation formula.
Bohr used the word "correspondence" in italics in lectures and writing before calling it a correspondence principle. He viewed this as a correspondence between quantum motion and radiation, not between classical and quantum theories. He writes in 1920 that there exists "a far-reaching correspondence between the various types of possible transitions between the stationary states on the one hand and the various harmonic components of the motion on the other hand."
Bohr's first article containing the definition of the correspondence principle was in 1923, in a summary paper entitled (in the English translation) "On the application of quantum theory to atomic structure". In his chapter II, "The process of radiation", he defines his correspondence principle as a condition connecting harmonic components of the electron moment to the possible occurrence of a radiative transition. In modern terms, this condition is a selection rule, saying that a given quantum jump is possible if and only if a particular type of motion exists in the corresponding classical model.
Following his definition of the correspondence principle, Bohr describes two applications. First he shows that the frequency of emitted radiation is related to an integral which can be well approximated by a sum when the quantum numbers inside the integral are large compared with their differences. Similarly he shows a relationship for the intensities of spectral lines and thus the rates at which quantum jumps occur.
These asymptotic relationships are expressed by Bohr as consequences of his general correspondence principle. However, historically each of these applications have been called "the correspondence principle".
The PhD dissertation of Hans Kramers working in Bohr's group in Copenhagen applied Bohr's correspondence principle to account for all of the known facts of the spectroscopic Stark effect, including some spectral components not known at the time of Kramers work.
Sommerfeld had been skeptical of the correspondence principle as it did not seem to be a consequence of a fundamental theory; Kramers' work convinced him that the principle had heuristic utility nevertheless. Other physicists picked up the concept, including work by John Van Vleck and by Kramers and Heisenberg on dispersion theory. The principle became a cornerstone of the semi-classical Bohr-Sommerfeld atomic theory;
Bohr's 1922 Nobel prize was partly awarded for his work with the correspondence principle.
Despite the successes, the physical theories based on the principle faced increasing challenges the early 1920s. Theoretical calculations by Van Vleck and by Kramers of the ionization potential of Helium disagreed significantly with experimental values. Bohr, Kramers, and John C. Slater responded with a new theoretical approach now called the BKS theory based on the correspondence principle but disavowing conservation of energy. Einstein and Wolfgang Pauli criticized the new approach, and the Bothe–Geiger coincidence experiment showed that energy was conserved in quantum collisions.
With the existing theories in conflict with observations, two new quantum mechanics concepts arose. First, Heisenberg's 1925 Umdeutung paper on matrix mechanics was inspired by the correspondence principle, although he did not cite Bohr. Further development in collaboration with Pascual Jordan and Max Born resulted in a mathematical model without connection to the principle. Second, Schrodinger's wave mechanics in the following year similarly did not use the principle. Both pictures were later shown to be equivalent and accurate enough to replace old quantum theory. These approaches have no atomic orbits: the correspondence is more of an analogy than a principle.
Dirac's correspondence
Paul Dirac developed significant portions of the new quantum theory in the second half of the 1920s. While he did not apply Bohr's correspondence principle, he developed a different, more formal classical–quantum correspondence. Dirac connected the structures of classical mechanics known as Poisson brackets to analogous structures of quantum mechanics known as commutators:
By this correspondence, now called canonical quantization, Dirac showed how the mathematical form of classical mechanics could be recast as a basis for the new mathematics of quantum mechanics.
Dirac developed these connections by studying the work of Heisenberg and Kramers on dispersion, work that was directly built on Bohr's correspondence principle; the Dirac approach provides a mathematically sound path towards Bohr's goal of a connection between classical and quantum mechanics. While Dirac did not call this correspondence a "principle", physics textbooks refer to his connections a "correspondence principle".
The classical limit of wave mechanics
The outstanding success of classical mechanics in the description of natural phenomena up to the 20th century means that quantum mechanics must do as well in similar circumstances.
One way to quantitatively define this concept is to require quantum mechanical theories to produce classical mechanics results as the quantum of action goes to zero, . This transition can be accomplished in two different ways.
First, the particle can be approximated by a wave packet, and the indefinite spread of the packet with time can be ignored. In 1927, Paul Ehrenfest proved his namesake theorem that showed that Newton's laws of motion hold on average in quantum mechanics: the quantum statistical expectation value of the position and momentum obey Newton's laws.
Second, the individual particle view can be replaced with a statistical mixture of classical particles with a density matching the quantum probability density. This approach led to the concept of semiclassical physics, beginning with the development of WKB approximation used in descriptions of quantum tunneling for example.
Modern view
While Bohr viewed "correspondence" as principle aiding his description of quantum phenomena, fundamental differences between the mathematical structure of quantum and of classical mechanics prevents correspondence in many cases. Rather than a principle, "there may be in some situations an approximate correspondence between classical and quantum concepts," physicist Asher Peres put it. Since quantum mechanics operates in a discrete space and classical mechanics in a continuous one, any correspondence will be necessarily fuzzy and elusive.
Introductory quantum mechanics textbooks suggest that quantum mechanics goes over to classical theory in the limit of high quantum numbers or in a limit where the Planck constant in the quantum formula is reduced to zero, . However such correspondence is not always possible. For example, classical systems can exhibit chaotic orbits which diverge but quantum states are unitary and maintain a fixed overlap.
Generalized correspondence principle
The term "generalized correspondence principle" has been used in the study of the history of science to mean the reduction of a new scientific theory to an earlier scientific theory in appropriate circumstances. This requires that the new theory explain all the phenomena under circumstances for which the preceding theory was known to be valid; it also means that new theory will retain large parts of the older theory. The generalized principle applies correspondence across aspects of a complete theory, not just a single formula as in the classical limit correspondence. For example, Albert Einstein in his 1905 work on relativity noted that classical mechanics relied on Galilean relativity while electromagnetism did not, and yet both work well. He produced a new theory that combined them in a away that reduced to these separate theories in approximations.
Ironically the singular failure of this "generalized correspondence principle" concept of scientific theories is the replacement of classical mechanics with quantum mechanics.
| Physical sciences | Quantum mechanics | Physics |
59868 | https://en.wikipedia.org/wiki/Interpreter%20%28computing%29 | Interpreter (computing) | In computer science, an interpreter is a computer program that directly executes instructions written in a programming or scripting language, without requiring them previously to have been compiled into a machine language program. An interpreter generally uses one of the following strategies for program execution:
Parse the source code and perform its behavior directly;
Translate source code into some efficient intermediate representation or object code and immediately execute that;
Explicitly execute stored precompiled bytecode made by a compiler and matched with the interpreter's virtual machine.
Early versions of Lisp programming language and minicomputer and microcomputer BASIC dialects would be examples of the first type. Perl, Raku, Python, MATLAB, and Ruby are examples of the second, while UCSD Pascal is an example of the third type. Source programs are compiled ahead of time and stored as machine independent code, which is then linked at run-time and executed by an interpreter and/or compiler (for JIT systems). Some systems, such as Smalltalk and contemporary versions of BASIC and Java, may also combine two and three types. Interpreters of various types have also been constructed for many languages traditionally associated with compilation, such as Algol, Fortran, Cobol, C and C++.
While interpretation and compilation are the two main means by which programming languages are implemented, they are not mutually exclusive, as most interpreting systems also perform some translation work, just like compilers. The terms "interpreted language" or "compiled language" signify that the canonical implementation of that language is an interpreter or a compiler, respectively. A high-level language is ideally an abstraction independent of particular implementations.
History
Interpreters were used as early as 1952 to ease programming within the limitations of computers at the time (e.g. a shortage of program storage space, or no native support for floating point numbers). Interpreters were also used to translate between low-level machine languages, allowing code to be written for machines that were still under construction and tested on computers that already existed. The first interpreted high-level language was Lisp. Lisp was first implemented by Steve Russell on an IBM 704 computer. Russell had read John McCarthy's paper, "Recursive Functions of Symbolic Expressions and Their Computation by Machine, Part I", and realized (to McCarthy's surprise) that the Lisp eval function could be implemented in machine code. The result was a working Lisp interpreter which could be used to run Lisp programs, or more properly, "evaluate Lisp expressions".
The development of editing interpreters was influenced by the need for interactive computing. In the 1960s, the introduction of time-sharing systems allowed multiple users to access a computer simultaneously, and editing interpreters became essential for managing and modifying code in real-time. The first editing interpreters were likely developed for mainframe computers, where they were used to create and modify programs on the fly. One of the earliest examples of an editing interpreter is the EDT (Editor and Debugger for the TECO) system, which was developed in the late 1960s for the PDP-1 computer. EDT allowed users to edit and debug programs using a combination of commands and macros, paving the way for modern text editors and interactive development environments.
General operation
An interpreter usually consists of a set of known commands it can execute, and a list of these commands in the order a programmer wishes to execute them. Each command (also known as an Instruction) contains the data the programmer wants to mutate, and information on how to mutate the data. For example, an interpreter might read ADD Books, 5 and interpret it as a request to add five to the Books variable.
Interpreters have a wide variety of instructions which are specialized to perform different tasks, but you will commonly find interpreter instructions for basic mathematical operations, branching, and memory management, making most interpreters Turing complete. Many interpreters are also closely integrated with a garbage collector and debugger.
Compilers versus interpreters
Programs written in a high-level language are either directly executed by some kind of interpreter or converted into machine code by a compiler (and assembler and linker) for the CPU to execute.
While compilers (and assemblers) generally produce machine code directly executable by computer hardware, they can often (optionally) produce an intermediate form called object code. This is basically the same machine specific code but augmented with a symbol table with names and tags to make executable blocks (or modules) identifiable and relocatable. Compiled programs will typically use building blocks (functions) kept in a library of such object code modules. A linker is used to combine (pre-made) library files with the object file(s) of the application to form a single executable file. The object files that are used to generate an executable file are thus often produced at different times, and sometimes even by different languages (capable of generating the same object format).
A simple interpreter written in a low-level language (e.g. assembly) may have similar machine code blocks implementing functions of the high-level language stored, and executed when a function's entry in a look up table points to that code. However, an interpreter written in a high-level language typically uses another approach, such as generating and then walking a parse tree, or by generating and executing intermediate software-defined instructions, or both.
Thus, both compilers and interpreters generally turn source code (text files) into tokens, both may (or may not) generate a parse tree, and both may generate immediate instructions (for a stack machine, quadruple code, or by other means). The basic difference is that a compiler system, including a (built in or separate) linker, generates a stand-alone machine code program, while an interpreter system instead performs the actions described by the high-level program.
A compiler can thus make almost all the conversions from source code semantics to the machine level once and for all (i.e. until the program has to be changed) while an interpreter has to do some of this conversion work every time a statement or function is executed. However, in an efficient interpreter, much of the translation work (including analysis of types, and similar) is factored out and done only the first time a program, module, function, or even statement, is run, thus quite akin to how a compiler works. However, a compiled program still runs much faster, under most circumstances, in part because compilers are designed to optimize code, and may be given ample time for this. This is especially true for simpler high-level languages without (many) dynamic data structures, checks, or type checking.
In traditional compilation, the executable output of the linkers (.exe files or .dll files or a library, see picture) is typically relocatable when run under a general operating system, much like the object code modules are but with the difference that this relocation is done dynamically at run time, i.e. when the program is loaded for execution. On the other hand, compiled and linked programs for small embedded systems are typically statically allocated, often hard coded in a NOR flash memory, as there is often no secondary storage and no operating system in this sense.
Historically, most interpreter systems have had a self-contained editor built in. This is becoming more common also for compilers (then often called an IDE), although some programmers prefer to use an editor of their choice and run the compiler, linker and other tools manually. Historically, compilers predate interpreters because hardware at that time could not support both the interpreter and interpreted code and the typical batch environment of the time limited the advantages of interpretation.
Development cycle
During the software development cycle, programmers make frequent changes to source code. When using a compiler, each time a change is made to the source code, they must wait for the compiler to translate the altered source files and link all of the binary code files together before the program can be executed. The larger the program, the longer the wait. By contrast, a programmer using an interpreter does a lot less waiting, as the interpreter usually just needs to translate the code being worked on to an intermediate representation (or not translate it at all), thus requiring much less time before the changes can be tested. Effects are evident upon saving the source code and reloading the program. Compiled code is generally less readily debugged as editing, compiling, and linking are sequential processes that have to be conducted in the proper sequence with a proper set of commands. For this reason, many compilers also have an executive aid, known as a Makefile and program. The Makefile lists compiler and linker command lines and program source code files, but might take a simple command line menu input (e.g. "Make 3") which selects the third group (set) of instructions then issues the commands to the compiler, and linker feeding the specified source code files.
Distribution
A compiler converts source code into binary instruction for a specific processor's architecture, thus making it less portable. This conversion is made just once, on the developer's environment, and after that the same binary can be distributed to the user's machines where it can be executed without further translation. A cross compiler can generate binary code for the user machine even if it has a different processor than the machine where the code is compiled.
An interpreted program can be distributed as source code. It needs to be translated in each final machine, which takes more time but makes the program distribution independent of the machine's architecture. However, the portability of interpreted source code is dependent on the target machine actually having a suitable interpreter. If the interpreter needs to be supplied along with the source, the overall installation process is more complex than delivery of a monolithic executable, since the interpreter itself is part of what needs to be installed.
The fact that interpreted code can easily be read and copied by humans can be of concern from the point of view of copyright. However, various systems of encryption and obfuscation exist. Delivery of intermediate code, such as bytecode, has a similar effect to obfuscation, but bytecode could be decoded with a decompiler or disassembler.
Efficiency
The main disadvantage of interpreters is that an interpreted program typically runs more slowly than if it had been compiled. The difference in speeds could be tiny or great; often an order of magnitude and sometimes more. It generally takes longer to run a program under an interpreter than to run the compiled code but it can take less time to interpret it than the total time required to compile and run it. This is especially important when prototyping and testing code when an edit-interpret-debug cycle can often be much shorter than an edit-compile-run-debug cycle.
Interpreting code is slower than running the compiled code because the interpreter must analyze each statement in the program each time it is executed and then perform the desired action, whereas the compiled code just performs the action within a fixed context determined by the compilation. This run-time analysis is known as "interpretive overhead". Access to variables is also slower in an interpreter because the mapping of identifiers to storage locations must be done repeatedly at run-time rather than at compile time.
There are various compromises between the development speed when using an interpreter and the execution speed when using a compiler. Some systems (such as some Lisps) allow interpreted and compiled code to call each other and to share variables. This means that once a routine has been tested and debugged under the interpreter it can be compiled and thus benefit from faster execution while other routines are being developed. Many interpreters do not execute the source code as it stands but convert it into some more compact internal form. Many BASIC interpreters replace keywords with single byte tokens which can be used to find the instruction in a jump table. A few interpreters, such as the PBASIC interpreter, achieve even higher levels of program compaction by using a bit-oriented rather than a byte-oriented program memory structure, where commands tokens occupy perhaps 5 bits, nominally "16-bit" constants are stored in a variable-length code requiring 3, 6, 10, or 18 bits, and address operands include a "bit offset". Many BASIC interpreters can store and read back their own tokenized internal representation.
An interpreter might well use the same lexical analyzer and parser as the compiler and then interpret the resulting abstract syntax tree. Example data type definitions for the latter, and a toy interpreter for syntax trees obtained from C expressions are shown in the box.
Regression
Interpretation cannot be used as the sole method of execution: even though an interpreter can itself be interpreted and so on, a directly executed program is needed somewhere at the bottom of the stack because the code being interpreted is not, by definition, the same as the machine code that the CPU can execute.
Variations
Bytecode interpreters
There is a spectrum of possibilities between interpreting and compiling, depending on the amount of analysis performed before the program is executed. For example, Emacs Lisp is compiled to bytecode, which is a highly compressed and optimized representation of the Lisp source, but is not machine code (and therefore not tied to any particular hardware). This "compiled" code is then interpreted by a bytecode interpreter (itself written in C). The compiled code in this case is machine code for a virtual machine, which is implemented not in hardware, but in the bytecode interpreter. Such compiling interpreters are sometimes also called compreters. In a bytecode interpreter each instruction starts with a byte, and therefore bytecode interpreters have up to 256 instructions, although not all may be used. Some bytecodes may take multiple bytes, and may be arbitrarily complicated.
Control tables - that do not necessarily ever need to pass through a compiling phase - dictate appropriate algorithmic control flow via customized interpreters in similar fashion to bytecode interpreters.
Threaded code interpreters
Threaded code interpreters are similar to bytecode interpreters but instead of bytes they use pointers. Each "instruction" is a word that points to a function or an instruction sequence, possibly followed by a parameter. The threaded code interpreter either loops fetching instructions and calling the functions they point to, or fetches the first instruction and jumps to it, and every instruction sequence ends with a fetch and jump to the next instruction. Unlike bytecode there is no effective limit on the number of different instructions other than available memory and address space. The classic example of threaded code is the Forth code used in Open Firmware systems: the source language is compiled into "F code" (a bytecode), which is then interpreted by a virtual machine.
Abstract syntax tree interpreters
In the spectrum between interpreting and compiling, another approach is to transform the source code into an optimized abstract syntax tree (AST), then execute the program following this tree structure, or use it to generate native code just-in-time. In this approach, each sentence needs to be parsed just once. As an advantage over bytecode, the AST keeps the global program structure and relations between statements (which is lost in a bytecode representation), and when compressed provides a more compact representation. Thus, using AST has been proposed as a better intermediate format for just-in-time compilers than bytecode. Also, it allows the system to perform better analysis during runtime.
However, for interpreters, an AST causes more overhead than a bytecode interpreter, because of nodes related to syntax performing no useful work, of a less sequential representation (requiring traversal of more pointers) and of overhead visiting the tree.
Just-in-time compilation
Further blurring the distinction between interpreters, bytecode interpreters and compilation is just-in-time (JIT) compilation, a technique in which the intermediate representation is compiled to native machine code at runtime. This confers the efficiency of running native code, at the cost of startup time and increased memory use when the bytecode or AST is first compiled. The earliest published JIT compiler is generally attributed to work on LISP by John McCarthy in 1960. Adaptive optimization is a complementary technique in which the interpreter profiles the running program and compiles its most frequently executed parts into native code. The latter technique is a few decades old, appearing in languages such as Smalltalk in the 1980s.
Just-in-time compilation has gained mainstream attention amongst language implementers in recent years, with Java, the .NET Framework, most modern JavaScript implementations, and Matlab now including JIT compilers.
Template Interpreter
Making the distinction between compilers and interpreters yet again even more vague is a special interpreter design known as a template interpreter. Rather than implement the execution of code by virtue of a large switch statement containing every possible bytecode, while operating on a software stack or a tree walk, a template interpreter maintains a large array of bytecode (or any efficient intermediate representation) mapped directly to corresponding native machine instructions that can be executed on the host hardware as key value pairs (or in more efficient designs, direct addresses to the native instructions), known as a "Template". When the particular code segment is executed the interpreter simply loads or jumps to the opcode mapping in the template and directly runs it on the hardware. Due to its design, the template interpreter very strongly resembles a just-in-time compiler rather than a traditional interpreter, however it is technically not a JIT due to the fact that it merely translates code from the language into native calls one opcode at a time rather than creating optimized sequences of CPU executable instructions from the entire code segment. Due to the interpreter's simple design of simply passing calls directly to the hardware rather than implementing them directly, it is much faster than every other type, even bytecode interpreters, and to an extent less prone to bugs, but as a tradeoff is more difficult to maintain due to the interpreter having to support translation to multiple different architectures instead of a platform independent virtual machine/stack. To date, the only template interpreter implementations of widely known languages to exist are the interpreter within Java's official reference implementation, the Sun HotSpot Java Virtual Machine, and the Ignition Interpreter in the Google V8 javascript execution engine.
Self-interpreter
A self-interpreter is a programming language interpreter written in a programming language which can interpret itself; an example is a BASIC interpreter written in BASIC. Self-interpreters are related to self-hosting compilers.
If no compiler exists for the language to be interpreted, creating a self-interpreter requires the implementation of the language in a host language (which may be another programming language or assembler). By having a first interpreter such as this, the system is bootstrapped and new versions of the interpreter can be developed in the language itself. It was in this way that Donald Knuth developed the TANGLE interpreter for the language WEB of the de-facto standard TeX typesetting system.
Defining a computer language is usually done in relation to an abstract machine (so-called operational semantics) or as a mathematical function (denotational semantics). A language may also be defined by an interpreter in which the semantics of the host language is given. The definition of a language by a self-interpreter is not well-founded (it cannot define a language), but a self-interpreter tells a reader about the expressiveness and elegance of a language. It also enables the interpreter to interpret its source code, the first step towards reflective interpreting.
An important design dimension in the implementation of a self-interpreter is whether a feature of the interpreted language is implemented with the same feature in the interpreter's host language. An example is whether a closure in a Lisp-like language is implemented using closures in the interpreter language or implemented "manually" with a data structure explicitly storing the environment. The more features implemented by the same feature in the host language, the less control the programmer of the interpreter has; for example, a different behavior for dealing with number overflows cannot be realized if the arithmetic operations are delegated to corresponding operations in the host language.
Some languages such as Lisp and Prolog have elegant self-interpreters. Much research on self-interpreters (particularly reflective interpreters) has been conducted in the Scheme programming language, a dialect of Lisp. In general, however, any Turing-complete language allows writing of its own interpreter. Lisp is such a language, because Lisp programs are lists of symbols and other lists. XSLT is such a language, because XSLT programs are written in XML. A sub-domain of metaprogramming is the writing of domain-specific languages (DSLs).
Clive Gifford introduced a measure quality of self-interpreter (the eigenratio), the limit of the ratio between computer time spent running a stack of N self-interpreters and time spent to run a stack of self-interpreters as N goes to infinity. This value does not depend on the program being run.
The book Structure and Interpretation of Computer Programs presents examples of meta-circular interpretation for Scheme and its dialects. Other examples of languages with a self-interpreter are Forth and Pascal.
Microcode
Microcode is a very commonly used technique "that imposes an interpreter between the hardware and the architectural level of a computer". As such, the microcode is a layer of hardware-level instructions that implement higher-level machine code instructions or internal state machine sequencing in many digital processing elements. Microcode is used in general-purpose central processing units, as well as in more specialized processors such as microcontrollers, digital signal processors, channel controllers, disk controllers, network interface controllers, network processors, graphics processing units, and in other hardware.
Microcode typically resides in special high-speed memory and translates machine instructions, state machine data or other input into sequences of detailed circuit-level operations. It separates the machine instructions from the underlying electronics so that instructions can be designed and altered more freely. It also facilitates the building of complex multi-step instructions, while reducing the complexity of computer circuits. Writing microcode is often called microprogramming and the microcode in a particular processor implementation is sometimes called a microprogram.
More extensive microcoding allows small and simple microarchitectures to emulate more powerful architectures with wider word length, more execution units and so on, which is a relatively simple way to achieve software compatibility between different products in a processor family.
Computer processor
Even a non microcoding computer processor itself can be considered to be a parsing immediate execution interpreter that is written in a general purpose hardware description language such as VHDL to create a system that parses the machine code instructions and immediately executes them.
Understanding Interpreter Performance
Interpreters, such as those written in Java, Perl, and Tcl, are now necessary for a wide range of computational tasks, including binary emulation and internet applications. Interpreter performance is still a worry despite their adaptability, particularly on systems with limited hardware resources. Advanced instrumentation and tracing approaches provide insights into interpreter implementations and processor resource utilization during execution through evaluations of interpreters tailored for the MIPS instruction set and programming languages such as Tcl, Perl, and Java. Performance characteristics are influenced by interpreter complexity, as demonstrated by comparisons with compiled code. It is clear that interpreter performance is more dependent on the nuances and resource needs of the interpreter than it is on the particular application that is being interpreted.
Applications
Interpreters are frequently used to execute command languages, and glue languages since each operator executed in command language is usually an invocation of a complex routine such as an editor or compiler.
Self-modifying code can easily be implemented in an interpreted language. This relates to the origins of interpretation in Lisp and artificial intelligence research.
Virtualization. Machine code intended for a hardware architecture can be run using a virtual machine. This is often used when the intended architecture is unavailable, or among other uses, for running multiple copies.
Sandboxing: While some types of sandboxes rely on operating system protections, an interpreter or virtual machine is often used. The actual hardware architecture and the originally intended hardware architecture may or may not be the same. This may seem pointless, except that sandboxes are not compelled to actually execute all the instructions the source code it is processing. In particular, it can refuse to execute code that violates any security constraints it is operating under.
Emulators for running computer software written for obsolete and unavailable hardware on more modern equipment.
| Technology | Software development: General | null |
59874 | https://en.wikipedia.org/wiki/Schr%C3%B6dinger%20equation | Schrödinger equation | The Schrödinger equation is a partial differential equation that governs the wave function of a non-relativistic quantum-mechanical system. Its discovery was a significant landmark in the development of quantum mechanics. It is named after Erwin Schrödinger, an Austrian physicist, who postulated the equation in 1925 and published it in 1926, forming the basis for the work that resulted in his Nobel Prize in Physics in 1933.
Conceptually, the Schrödinger equation is the quantum counterpart of Newton's second law in classical mechanics. Given a set of known initial conditions, Newton's second law makes a mathematical prediction as to what path a given physical system will take over time. The Schrödinger equation gives the evolution over time of the wave function, the quantum-mechanical characterization of an isolated physical system. The equation was postulated by Schrödinger based on a postulate of Louis de Broglie that all matter has an associated matter wave. The equation predicted bound states of the atom in agreement with experimental observations.
The Schrödinger equation is not the only way to study quantum mechanical systems and make predictions. Other formulations of quantum mechanics include matrix mechanics, introduced by Werner Heisenberg, and the path integral formulation, developed chiefly by Richard Feynman. When these approaches are compared, the use of the Schrödinger equation is sometimes called "wave mechanics". The Klein-Gordon equation is a wave equation which is the relativistic version of the Schrödinger equation. The Schrödinger equation is nonrelativistic because it contains a first derivative in time and a second derivative in space, and therefore space & time are not on equal footing.
Paul Dirac incorporated special relativity and quantum mechanics into a single formulation that simplifies to the Schrödinger equation in the non-relativistic limit. This is the Dirac equation, which contains a single derivative in both space and time. The second-derivative PDE of the Klein-Gordon equation led to a problem with probability density even though it was a relativistic wave equation. The probability density could be negative, which is physically unviable. This was fixed by Dirac by taking the so-called square-root of the Klein-Gordon operator and in turn introducing Dirac matrices. In a modern context, the Klein-Gordon equation describes spin-less particles, while the Dirac equation describes spin-1/2 particles.
Definition
Preliminaries
Introductory courses on physics or chemistry typically introduce the Schrödinger equation in a way that can be appreciated knowing only the concepts and notations of basic calculus, particularly derivatives with respect to space and time. A special case of the Schrödinger equation that admits a statement in those terms is the position-space Schrödinger equation for a single nonrelativistic particle in one dimension:
Here, is a wave function, a function that assigns a complex number to each point at each time . The parameter is the mass of the particle, and is the potential that represents the environment in which the particle exists. The constant is the imaginary unit, and is the reduced Planck constant, which has units of action (energy multiplied by time).
Broadening beyond this simple case, the mathematical formulation of quantum mechanics developed by Paul Dirac, David Hilbert, John von Neumann, and Hermann Weyl defines the state of a quantum mechanical system to be a vector belonging to a separable complex Hilbert space . This vector is postulated to be normalized under the Hilbert space's inner product, that is, in Dirac notation it obeys . The exact nature of this Hilbert space is dependent on the system – for example, for describing position and momentum the Hilbert space is the space of square-integrable functions , while the Hilbert space for the spin of a single proton is the two-dimensional complex vector space with the usual inner product.
Physical quantities of interest – position, momentum, energy, spin – are represented by observables, which are self-adjoint operators acting on the Hilbert space. A wave function can be an eigenvector of an observable, in which case it is called an eigenstate, and the associated eigenvalue corresponds to the value of the observable in that eigenstate. More generally, a quantum state will be a linear combination of the eigenstates, known as a quantum superposition. When an observable is measured, the result will be one of its eigenvalues with probability given by the Born rule: in the simplest case the eigenvalue is non-degenerate and the probability is given by , where is its associated eigenvector. More generally, the eigenvalue is degenerate and the probability is given by , where is the projector onto its associated eigenspace.
A momentum eigenstate would be a perfectly monochromatic wave of infinite extent, which is not square-integrable. Likewise a position eigenstate would be a Dirac delta distribution, not square-integrable and technically not a function at all. Consequently, neither can belong to the particle's Hilbert space. Physicists sometimes regard these eigenstates, composed of elements outside the Hilbert space, as "generalized eigenvectors". These are used for calculational convenience and do not represent physical states. Thus, a position-space wave function as used above can be written as the inner product of a time-dependent state vector with unphysical but convenient "position eigenstates" :
Time-dependent equation
The form of the Schrödinger equation depends on the physical situation. The most general form is the time-dependent Schrödinger equation, which gives a description of a system evolving with time:
where is time, is the state vector of the quantum system ( being the Greek letter psi), and is an observable, the Hamiltonian operator.
The term "Schrödinger equation" can refer to both the general equation, or the specific nonrelativistic version. The general equation is indeed quite general, used throughout quantum mechanics, for everything from the Dirac equation to quantum field theory, by plugging in diverse expressions for the Hamiltonian. The specific nonrelativistic version is an approximation that yields accurate results in many situations, but only to a certain extent (see relativistic quantum mechanics and relativistic quantum field theory).
To apply the Schrödinger equation, write down the Hamiltonian for the system, accounting for the kinetic and potential energies of the particles constituting the system, then insert it into the Schrödinger equation. The resulting partial differential equation is solved for the wave function, which contains information about the system. In practice, the square of the absolute value of the wave function at each point is taken to define a probability density function. For example, given a wave function in position space as above, we have
Time-independent equation
The time-dependent Schrödinger equation described above predicts that wave functions can form standing waves, called stationary states. These states are particularly important as their individual study later simplifies the task of solving the time-dependent Schrödinger equation for any state. Stationary states can also be described by a simpler form of the Schrödinger equation, the time-independent Schrödinger equation.
where is the energy of the system. This is only used when the Hamiltonian itself is not dependent on time explicitly. However, even in this case the total wave function is dependent on time as explained in the section on linearity below. In the language of linear algebra, this equation is an eigenvalue equation. Therefore, the wave function is an eigenfunction of the Hamiltonian operator with corresponding eigenvalue(s) .
Properties
Linearity
The Schrödinger equation is a linear differential equation, meaning that if two state vectors and are solutions, then so is any linear combination
of the two state vectors where and are any complex numbers. Moreover, the sum can be extended for any number of state vectors. This property allows superpositions of quantum states to be solutions of the Schrödinger equation. Even more generally, it holds that a general solution to the Schrödinger equation can be found by taking a weighted sum over a basis of states. A choice often employed is the basis of energy eigenstates, which are solutions of the time-independent Schrödinger equation. In this basis, a time-dependent state vector can be written as the linear combination
where are complex numbers and the vectors are solutions of the time-independent equation .
Unitarity
Holding the Hamiltonian constant, the Schrödinger equation has the solution
The operator is known as the time-evolution operator, and it is unitary: it preserves the inner product between vectors in the Hilbert space. Unitarity is a general feature of time evolution under the Schrödinger equation. If the initial state is , then the state at a later time will be given by
for some unitary operator . Conversely, suppose that is a continuous family of unitary operators parameterized by . Without loss of generality, the parameterization can be chosen so that is the identity operator and that for any . Then depends upon the parameter in such a way that
for some self-adjoint operator , called the generator of the family . A Hamiltonian is just such a generator (up to the factor of the Planck constant that would be set to 1 in natural units).
To see that the generator is Hermitian, note that with , we have
so is unitary only if, to first order, its derivative is Hermitian.
Changes of basis
The Schrödinger equation is often presented using quantities varying as functions of position, but as a vector-operator equation it has a valid representation in any arbitrary complete basis of kets in Hilbert space. As mentioned above, "bases" that lie outside the physical Hilbert space are also employed for calculational purposes. This is illustrated by the position-space and momentum-space Schrödinger equations for a nonrelativistic, spinless particle. The Hilbert space for such a particle is the space of complex square-integrable functions on three-dimensional Euclidean space, and its Hamiltonian is the sum of a kinetic-energy term that is quadratic in the momentum operator and a potential-energy term:
Writing for a three-dimensional position vector and for a three-dimensional momentum vector, the position-space Schrödinger equation is
The momentum-space counterpart involves the Fourier transforms of the wave function and the potential:
The functions and are derived from by
where and do not belong to the Hilbert space itself, but have well-defined inner products with all elements of that space.
When restricted from three dimensions to one, the position-space equation is just the first form of the Schrödinger equation given above. The relation between position and momentum in quantum mechanics can be appreciated in a single dimension. In canonical quantization, the classical variables and are promoted to self-adjoint operators and that satisfy the canonical commutation relation
This implies that
so the action of the momentum operator in the position-space representation is . Thus, becomes a second derivative, and in three dimensions, the second derivative becomes the Laplacian .
The canonical commutation relation also implies that the position and momentum operators are Fourier conjugates of each other. Consequently, functions originally defined in terms of their position dependence can be converted to functions of momentum using the Fourier transform. In solid-state physics, the Schrödinger equation is often written for functions of momentum, as Bloch's theorem ensures the periodic crystal lattice potential couples with for only discrete reciprocal lattice vectors . This makes it convenient to solve the momentum-space Schrödinger equation at each point in the Brillouin zone independently of the other points in the Brillouin zone.
Probability current
The Schrödinger equation is consistent with local probability conservation. It also ensures that a normalized wavefunction remains normalized after time evolution. In matrix mechanics, this means that the time evolution operator is a unitary operator. In contrast to, for example, the Klein Gordon equation, although a redefined inner product of a wavefunction can be time independent, the total volume integral of modulus square of the wavefunction need not be time independent.
The continuity equation for probability in non relativistic quantum mechanics is stated as:
where
is the probability current or probability flux (flow per unit area).
If the wavefunction is represented as where is a real function which represents the complex phase of the wavefunction, then the probability flux is calculated as:Hence, the spatial variation of the phase of a wavefunction is said to characterize the probability flux of the wavefunction. Although the term appears to play the role of velocity, it does not represent velocity at a point since simultaneous measurement of position and velocity violates uncertainty principle.
Separation of variables
If the Hamiltonian is not an explicit function of time, Schrödinger's equation reads:
The operator on the left side depends only on time; the one on the right side depends only on space.
Solving the equation by separation of variables means seeking a solution of the form of a product of spatial and temporal parts
where is a function of all the spatial coordinate(s) of the particle(s) constituting the system only, and is a function of time only. Substituting this expression for into the time dependent left hand side shows that is a phase factor:
A solution of this type is called stationary, since the only time dependence is a phase factor that cancels when the probability density is calculated via the Born rule.
The spatial part of the full wave function solves:
where the energy appears in the phase factor.
This generalizes to any number of particles in any number of dimensions (in a time-independent potential): the standing wave solutions of the time-independent equation are the states with definite energy, instead of a probability distribution of different energies. In physics, these standing waves are called "stationary states" or "energy eigenstates"; in chemistry they are called "atomic orbitals" or "molecular orbitals". Superpositions of energy eigenstates change their properties according to the relative phases between the energy levels. The energy eigenstates form a basis: any wave function may be written as a sum over the discrete energy states or an integral over continuous energy states, or more generally as an integral over a measure. This is the spectral theorem in mathematics, and in a finite-dimensional state space it is just a statement of the completeness of the eigenvectors of a Hermitian matrix.
Separation of variables can also be a useful method for the time-independent Schrödinger equation. For example, depending on the symmetry of the problem, the Cartesian axes might be separated,
or radial and angular coordinates might be separated:
Examples
Particle in a box
The particle in a one-dimensional potential energy box is the most mathematically simple example where restraints lead to the quantization of energy levels. The box is defined as having zero potential energy inside a certain region and infinite potential energy outside. For the one-dimensional case in the direction, the time-independent Schrödinger equation may be written
With the differential operator defined by
the previous equation is evocative of the classic kinetic energy analogue,
with state in this case having energy coincident with the kinetic energy of the particle.
The general solutions of the Schrödinger equation for the particle in a box are
or, from Euler's formula,
The infinite potential walls of the box determine the values of and at and where must be zero. Thus, at ,
and . At ,
in which cannot be zero as this would conflict with the postulate that has norm 1. Therefore, since , must be an integer multiple of ,
This constraint on implies a constraint on the energy levels, yielding
A finite potential well is the generalization of the infinite potential well problem to potential wells having finite depth. The finite potential well problem is mathematically more complicated than the infinite particle-in-a-box problem as the wave function is not pinned to zero at the walls of the well. Instead, the wave function must satisfy more complicated mathematical boundary conditions as it is nonzero in regions outside the well. Another related problem is that of the rectangular potential barrier, which furnishes a model for the quantum tunneling effect that plays an important role in the performance of modern technologies such as flash memory and scanning tunneling microscopy.
Harmonic oscillator
The Schrödinger equation for this situation is
where is the displacement and the angular frequency. Furthermore, it can be used to describe approximately a wide variety of other systems, including vibrating atoms, molecules, and atoms or ions in lattices, and approximating other potentials near equilibrium points. It is also the basis of perturbation methods in quantum mechanics.
The solutions in position space are
where , and the functions are the Hermite polynomials of order . The solution set may be generated by
The eigenvalues are
The case is called the ground state, its energy is called the zero-point energy, and the wave function is a Gaussian.
The harmonic oscillator, like the particle in a box, illustrates the generic feature of the Schrödinger equation that the energies of bound eigenstates are discretized.
Hydrogen atom
The Schrödinger equation for the electron in a hydrogen atom (or a hydrogen-like atom) is
where is the electron charge, is the position of the electron relative to the nucleus, is the magnitude of the relative position, the potential term is due to the Coulomb interaction, wherein is the permittivity of free space and
is the 2-body reduced mass of the hydrogen nucleus (just a proton) of mass and the electron of mass . The negative sign arises in the potential term since the proton and electron are oppositely charged. The reduced mass in place of the electron mass is used since the electron and proton together orbit each other about a common center of mass, and constitute a two-body problem to solve. The motion of the electron is of principal interest here, so the equivalent one-body problem is the motion of the electron using the reduced mass.
The Schrödinger equation for a hydrogen atom can be solved by separation of variables. In this case, spherical polar coordinates are the most convenient. Thus,
where are radial functions and are spherical harmonics of degree and order . This is the only atom for which the Schrödinger equation has been solved for exactly. Multi-electron atoms require approximate methods. The family of solutions are:
where
is the Bohr radius,
are the generalized Laguerre polynomials of degree ,
are the principal, azimuthal, and magnetic quantum numbers respectively, which take the values
Approximate solutions
It is typically not possible to solve the Schrödinger equation exactly for situations of physical interest. Accordingly, approximate solutions are obtained using techniques like variational methods and WKB approximation. It is also common to treat a problem of interest as a small modification to a problem that can be solved exactly, a method known as perturbation theory.
Semiclassical limit
One simple way to compare classical to quantum mechanics is to consider the time-evolution of the expected position and expected momentum, which can then be compared to the time-evolution of the ordinary position and momentum in classical mechanics. The quantum expectation values satisfy the Ehrenfest theorem. For a one-dimensional quantum particle moving in a potential , the Ehrenfest theorem says
Although the first of these equations is consistent with the classical behavior, the second is not: If the pair were to satisfy Newton's second law, the right-hand side of the second equation would have to be
which is typically not the same as . For a general , therefore, quantum mechanics can lead to predictions where expectation values do not mimic the classical behavior. In the case of the quantum harmonic oscillator, however, is linear and this distinction disappears, so that in this very special case, the expected position and expected momentum do exactly follow the classical trajectories.
For general systems, the best we can hope for is that the expected position and momentum will approximately follow the classical trajectories. If the wave function is highly concentrated around a point , then and will be almost the same, since both will be approximately equal to . In that case, the expected position and expected momentum will remain very close to the classical trajectories, at least for as long as the wave function remains highly localized in position.
The Schrödinger equation in its general form
is closely related to the Hamilton–Jacobi equation (HJE)
where is the classical action and is the Hamiltonian function (not operator). Here the generalized coordinates for (used in the context of the HJE) can be set to the position in Cartesian coordinates as .
Substituting
where is the probability density, into the Schrödinger equation and then taking the limit in the resulting equation yield the Hamilton–Jacobi equation.
Density matrices
Wave functions are not always the most convenient way to describe quantum systems and their behavior. When the preparation of a system is only imperfectly known, or when the system under investigation is a part of a larger whole, density matrices may be used instead. A density matrix is a positive semi-definite operator whose trace is equal to 1. (The term "density operator" is also used, particularly when the underlying Hilbert space is infinite-dimensional.) The set of all density matrices is convex, and the extreme points are the operators that project onto vectors in the Hilbert space. These are the density-matrix representations of wave functions; in Dirac notation, they are written
The density-matrix analogue of the Schrödinger equation for wave functions is
where the brackets denote a commutator. This is variously known as the von Neumann equation, the Liouville–von Neumann equation, or just the Schrödinger equation for density matrices. If the Hamiltonian is time-independent, this equation can be easily solved to yield
More generally, if the unitary operator describes wave function evolution over some time interval, then the time evolution of a density matrix over that same interval is given by
Unitary evolution of a density matrix conserves its von Neumann entropy.
Relativistic quantum physics and quantum field theory
The one-particle Schrödinger equation described above is valid essentially in the nonrelativistic domain. For one reason, it is essentially invariant under Galilean transformations, which form the symmetry group of Newtonian dynamics. Moreover, processes that change particle number are natural in relativity, and so an equation for one particle (or any fixed number thereof) can only be of limited use. A more general form of the Schrödinger equation that also applies in relativistic situations can be formulated within quantum field theory (QFT), a framework that allows the combination of quantum mechanics with special relativity. The region in which both simultaneously apply may be described by relativistic quantum mechanics. Such descriptions may use time evolution generated by a Hamiltonian operator, as in the Schrödinger functional method.
Klein–Gordon and Dirac equations
Attempts to combine quantum physics with special relativity began with building relativistic wave equations from the relativistic energy–momentum relation
instead of nonrelativistic energy equations. The Klein–Gordon equation and the Dirac equation are two such equations. The Klein–Gordon equation,
was the first such equation to be obtained, even before the nonrelativistic one-particle Schrödinger equation, and applies to massive spinless particles. Historically, Dirac obtained the Dirac equation by seeking a differential equation that would be first-order in both time and space, a desirable property for a relativistic theory. Taking the "square root" of the left-hand side of the Klein–Gordon equation in this way required factorizing it into a product of two operators, which Dirac wrote using 4 × 4 matrices . Consequently, the wave function also became a four-component function, governed by the Dirac equation that, in free space, read
This has again the form of the Schrödinger equation, with the time derivative of the wave function being given by a Hamiltonian operator acting upon the wave function. Including influences upon the particle requires modifying the Hamiltonian operator. For example, the Dirac Hamiltonian for a particle of mass and electric charge in an electromagnetic field (described by the electromagnetic potentials and ) is:
in which the and are the Dirac gamma matrices related to the spin of the particle. The Dirac equation is true for all particles, and the solutions to the equation are spinor fields with two components corresponding to the particle and the other two for the antiparticle.
For the Klein–Gordon equation, the general form of the Schrödinger equation is inconvenient to use, and in practice the Hamiltonian is not expressed in an analogous way to the Dirac Hamiltonian. The equations for relativistic quantum fields, of which the Klein–Gordon and Dirac equations are two examples, can be obtained in other ways, such as starting from a Lagrangian density and using the Euler–Lagrange equations for fields, or using the representation theory of the Lorentz group in which certain representations can be used to fix the equation for a free particle of given spin (and mass).
In general, the Hamiltonian to be substituted in the general Schrödinger equation is not just a function of the position and momentum operators (and possibly time), but also of spin matrices. Also, the solutions to a relativistic wave equation, for a massive particle of spin , are complex-valued spinor fields.
Fock space
As originally formulated, the Dirac equation is an equation for a single quantum particle, just like the single-particle Schrödinger equation with wave function This is of limited use in relativistic quantum mechanics, where particle number is not fixed. Heuristically, this complication can be motivated by noting that mass–energy equivalence implies material particles can be created from energy. A common way to address this in QFT is to introduce a Hilbert space where the basis states are labeled by particle number, a so-called Fock space. The Schrödinger equation can then be formulated for quantum states on this Hilbert space. However, because the Schrödinger equation picks out a preferred time axis, the Lorentz invariance of the theory is no longer manifest, and accordingly, the theory is often formulated in other ways.
History
Following Max Planck's quantization of light (see black-body radiation), Albert Einstein interpreted Planck's quanta to be photons, particles of light, and proposed that the energy of a photon is proportional to its frequency, one of the first signs of wave–particle duality. Since energy and momentum are related in the same way as frequency and wave number in special relativity, it followed that the momentum of a photon is inversely proportional to its wavelength , or proportional to its wave number :
where is the Planck constant and is the reduced Planck constant. Louis de Broglie hypothesized that this is true for all particles, even particles which have mass such as electrons. He showed that, assuming that the matter waves propagate along with their particle counterparts, electrons form standing waves, meaning that only certain discrete rotational frequencies about the nucleus of an atom are allowed.
These quantized orbits correspond to discrete energy levels, and de Broglie reproduced the Bohr model formula for the energy levels. The Bohr model was based on the assumed quantization of angular momentum according to
According to de Broglie, the electron is described by a wave, and a whole number of wavelengths must fit along the circumference of the electron's orbit:
This approach essentially confined the electron wave in one dimension, along a circular orbit of radius .
In 1921, prior to de Broglie, Arthur C. Lunn at the University of Chicago had used the same argument based on the completion of the relativistic energy–momentum 4-vector to derive what we now call the de Broglie relation. Unlike de Broglie, Lunn went on to formulate the differential equation now known as the Schrödinger equation and solve for its energy eigenvalues for the hydrogen atom; the paper was rejected by the Physical Review, according to Kamen.
Following up on de Broglie's ideas, physicist Peter Debye made an offhand comment that if particles behaved as waves, they should satisfy some sort of wave equation. Inspired by Debye's remark, Schrödinger decided to find a proper 3-dimensional wave equation for the electron. He was guided by William Rowan Hamilton's analogy between mechanics and optics, encoded in the observation that the zero-wavelength limit of optics resembles a mechanical system—the trajectories of light rays become sharp tracks that obey Fermat's principle, an analog of the principle of least action.
The equation he found is
By that time Arnold Sommerfeld had refined the Bohr model with relativistic corrections. Schrödinger used the relativistic energy–momentum relation to find what is now known as the Klein–Gordon equation in a Coulomb potential (in natural units):
He found the standing waves of this relativistic equation, but the relativistic corrections disagreed with Sommerfeld's formula. Discouraged, he put away his calculations and secluded himself with a mistress in a mountain cabin in December 1925.
While at the cabin, Schrödinger decided that his earlier nonrelativistic calculations were novel enough to publish and decided to leave off the problem of relativistic corrections for the future. Despite the difficulties in solving the differential equation for hydrogen (he had sought help from his friend the mathematician Hermann Weyl) Schrödinger showed that his nonrelativistic version of the wave equation produced the correct spectral energies of hydrogen in a paper published in 1926. Schrödinger computed the hydrogen spectral series by treating a hydrogen atom's electron as a wave , moving in a potential well , created by the proton. This computation accurately reproduced the energy levels of the Bohr model.
The Schrödinger equation details the behavior of but says nothing of its nature. Schrödinger tried to interpret the real part of as a charge density, and then revised this proposal, saying in his next paper that the modulus squared of is a charge density. This approach was, however, unsuccessful. In 1926, just a few days after this paper was published, Max Born successfully interpreted as the probability amplitude, whose modulus squared is equal to probability density. Later, Schrödinger himself explained this interpretation as follows:
Interpretation
The Schrödinger equation provides a way to calculate the wave function of a system and how it changes dynamically in time. However, the Schrödinger equation does not directly say what, exactly, the wave function is. The meaning of the Schrödinger equation and how the mathematical entities in it relate to physical reality depends upon the interpretation of quantum mechanics that one adopts.
In the views often grouped together as the Copenhagen interpretation, a system's wave function is a collection of statistical information about that system. The Schrödinger equation relates information about the system at one time to information about it at another. While the time-evolution process represented by the Schrödinger equation is continuous and deterministic, in that knowing the wave function at one instant is in principle sufficient to calculate it for all future times, wave functions can also change discontinuously and stochastically during a measurement. The wave function changes, according to this school of thought, because new information is available. The post-measurement wave function generally cannot be known prior to the measurement, but the probabilities for the different possibilities can be calculated using the Born rule. Other, more recent interpretations of quantum mechanics, such as relational quantum mechanics and QBism also give the Schrödinger equation a status of this sort.
Schrödinger himself suggested in 1952 that the different terms of a superposition evolving under the Schrödinger equation are "not alternatives but all really happen simultaneously". This has been interpreted as an early version of Everett's many-worlds interpretation. This interpretation, formulated independently in 1956, holds that all the possibilities described by quantum theory simultaneously occur in a multiverse composed of mostly independent parallel universes. This interpretation removes the axiom of wave function collapse, leaving only continuous evolution under the Schrödinger equation, and so all possible states of the measured system and the measuring apparatus, together with the observer, are present in a real physical quantum superposition. While the multiverse is deterministic, we perceive non-deterministic behavior governed by probabilities, because we do not observe the multiverse as a whole, but only one parallel universe at a time. Exactly how this is supposed to work has been the subject of much debate. Why we should assign probabilities at all to outcomes that are certain to occur in some worlds, and why should the probabilities be given by the Born rule? Several ways to answer these questions in the many-worlds framework have been proposed, but there is no consensus on whether they are successful.
Bohmian mechanics reformulates quantum mechanics to make it deterministic, at the price of adding a force due to a "quantum potential". It attributes to each physical system not only a wave function but in addition a real position that evolves deterministically under a nonlocal guiding equation. The evolution of a physical system is given at all times by the Schrödinger equation together with the guiding equation.
| Physical sciences | Quantum mechanics | null |
59881 | https://en.wikipedia.org/wiki/Ideal%20gas%20law | Ideal gas law | The ideal gas law, also called the general gas equation, is the equation of state of a hypothetical ideal gas. It is a good approximation of the behavior of many gases under many conditions, although it has several limitations. It was first stated by Benoît Paul Émile Clapeyron in 1834 as a combination of the empirical Boyle's law, Charles's law, Avogadro's law, and Gay-Lussac's law. The ideal gas law is often written in an empirical form:
where , and are the pressure, volume and temperature respectively; is the amount of substance; and is the ideal gas constant.
It can also be derived from the microscopic kinetic theory, as was achieved (apparently independently) by August Krönig in 1856 and Rudolf Clausius in 1857.
Equation
The state of an amount of gas is determined by its pressure, volume, and temperature. The modern form of the equation relates these simply in two main forms. The temperature used in the equation of state is an absolute temperature: the appropriate SI unit is the kelvin.
Common forms
The most frequently introduced forms are:where:
is the absolute pressure of the gas,
is the volume of the gas,
is the amount of substance of gas (also known as number of moles),
is the ideal, or universal, gas constant, equal to the product of the Boltzmann constant and the Avogadro constant,
is the Boltzmann constant,
is the Avogadro constant,
is the absolute temperature of the gas,
is the number of particles (usually atoms or molecules) of the gas.
In SI units, p is measured in pascals, V is measured in cubic metres, n is measured in moles, and T in kelvins (the Kelvin scale is a shifted Celsius scale, where 0 K = −273.15 °C, the lowest possible temperature). R has for value 8.314 J/(mol·K) = 1.989 ≈ 2 cal/(mol·K), or 0.0821 L⋅atm/(mol⋅K).
Molar form
How much gas is present could be specified by giving the mass instead of the chemical amount of gas. Therefore, an alternative form of the ideal gas law may be useful. The chemical amount, n (in moles), is equal to total mass of the gas (m) (in kilograms) divided by the molar mass, M (in kilograms per mole):
By replacing n with m/M and subsequently introducing density ρ = m/V, we get:
Defining the specific gas constant Rspecific as the ratio R/M,
This form of the ideal gas law is very useful because it links pressure, density, and temperature in a unique formula independent of the quantity of the considered gas. Alternatively, the law may be written in terms of the specific volume v, the reciprocal of density, as
It is common, especially in engineering and meteorological applications, to represent the specific gas constant by the symbol R. In such cases, the universal gas constant is usually given a different symbol such as or to distinguish it. In any case, the context and/or units of the gas constant should make it clear as to whether the universal or specific gas constant is being used.
Statistical mechanics
In statistical mechanics, the following molecular equation is derived from first principles
where is the absolute pressure of the gas, is the number density of the molecules (given by the ratio , in contrast to the previous formulation in which is the number of moles), is the absolute temperature, and is the Boltzmann constant relating temperature and energy, given by:
where is the Avogadro constant.
From this we notice that for a gas of mass , with an average particle mass of times the atomic mass constant, , (i.e., the mass is Da) the number of molecules will be given by
and since , we find that the ideal gas law can be rewritten as
In SI units, is measured in pascals, in cubic metres, in kelvins, and in SI units.
Combined gas law
Combining the laws of Charles, Boyle and Gay-Lussac gives the combined gas law, which takes the same functional form as the ideal gas law says that the number of moles is unspecified, and the ratio of to is simply taken as a constant:
where is the pressure of the gas, is the volume of the gas, is the absolute temperature of the gas, and is a constant. When comparing the same substance under two different sets of conditions, the law can be written as
Energy associated with a gas
According to the assumptions of the kinetic theory of ideal gases, one can consider that there are no intermolecular attractions between the molecules, or atoms, of an ideal gas. In other words, its potential energy is zero. Hence, all the energy possessed by the gas is the kinetic energy of the molecules, or atoms, of the gas.
This corresponds to the kinetic energy of n moles of a monoatomic gas having 3 degrees of freedom; x, y, z. The table here below gives this relationship for different amounts of a monoatomic gas.
Applications to thermodynamic processes
The table below essentially simplifies the ideal gas equation for a particular process, making the equation easier to solve using numerical methods.
A thermodynamic process is defined as a system that moves from state 1 to state 2, where the state number is denoted by a subscript. As shown in the first column of the table, basic thermodynamic processes are defined such that one of the gas properties (P, V, T, S, or H) is constant throughout the process.
For a given thermodynamic process, in order to specify the extent of a particular process, one of the properties ratios (which are listed under the column labeled "known ratio") must be specified (either directly or indirectly). Also, the property for which the ratio is known must be distinct from the property held constant in the previous column (otherwise the ratio would be unity, and not enough information would be available to simplify the gas law equation).
In the final three columns, the properties (p, V, or T) at state 2 can be calculated from the properties at state 1 using the equations listed.
a. In an isentropic process, system entropy (S) is constant. Under these conditions, p1V1γ = p2V2γ, where γ is defined as the heat capacity ratio, which is constant for a calorifically perfect gas. The value used for γ is typically 1.4 for diatomic gases like nitrogen (N2) and oxygen (O2), (and air, which is 99% diatomic). Also γ is typically 1.6 for mono atomic gases like the noble gases helium (He), and argon (Ar). In internal combustion engines γ varies between 1.35 and 1.15, depending on constitution gases and temperature.
b. In an isenthalpic process, system enthalpy (H) is constant. In the case of free expansion for an ideal gas, there are no molecular interactions, and the temperature remains constant. For real gasses, the molecules do interact via attraction or repulsion depending on temperature and pressure, and heating or cooling does occur. This is known as the Joule–Thomson effect. For reference, the Joule–Thomson coefficient μJT for air at room temperature and sea level is 0.22 °C/bar.
Deviations from ideal behavior of real gases
The equation of state given here (PV = nRT) applies only to an ideal gas, or as an approximation to a real gas that behaves sufficiently like an ideal gas. There are in fact many different forms of the equation of state. Since the ideal gas law neglects both molecular size and intermolecular attractions, it is most accurate for monatomic gases at high temperatures and low pressures. The neglect of molecular size becomes less important for lower densities, i.e. for larger volumes at lower pressures, because the average distance between adjacent molecules becomes much larger than the molecular size. The relative importance of intermolecular attractions diminishes with increasing thermal kinetic energy, i.e., with increasing temperatures. More detailed equations of state, such as the van der Waals equation, account for deviations from ideality caused by molecular size and intermolecular forces.
Derivations
Empirical
The empirical laws that led to the derivation of the ideal gas law were discovered with experiments that changed only 2 state variables of the gas and kept every other one constant.
All the possible gas laws that could have been discovered with this kind of setup are:
Boyle's law ()
Charles's law ()
Avogadro's law ()
Gay-Lussac's law ()
where P stands for pressure, V for volume, N for number of particles in the gas and T for temperature; where are constants in this context because of each equation requiring only the parameters explicitly noted in them changing.
To derive the ideal gas law one does not need to know all 6 formulas, one can just know 3 and with those derive the rest or just one more to be able to get the ideal gas law, which needs 4.
Since each formula only holds when only the state variables involved in said formula change while the others (which are a property of the gas but are not explicitly noted in said formula) remain constant, we cannot simply use algebra and directly combine them all. This is why: Boyle did his experiments while keeping N and T constant and this must be taken into account (in this same way, every experiment kept some parameter as constant and this must be taken into account for the derivation).
Keeping this in mind, to carry the derivation on correctly, one must imagine the gas being altered by one process at a time (as it was done in the experiments). The derivation using 4 formulas can look like this:
at first the gas has parameters
Say, starting to change only pressure and volume, according to Boyle's law (), then:
After this process, the gas has parameters
Using then equation () to change the number of particles in the gas and the temperature,
After this process, the gas has parameters
Using then equation () to change the pressure and the number of particles,
After this process, the gas has parameters
Using then Charles's law (equation 2) to change the volume and temperature of the gas,
After this process, the gas has parameters
Using simple algebra on equations (), (), () and () yields the result:
or where stands for the Boltzmann constant.
Another equivalent result, using the fact that , where n is the number of moles in the gas and R is the universal gas constant, is:
which is known as the ideal gas law.
If three of the six equations are known, it may be possible to derive the remaining three using the same method. However, because each formula has two variables, this is possible only for certain groups of three. For example, if you were to have equations (), () and () you would not be able to get any more because combining any two of them will only give you the third. However, if you had equations (), () and () you would be able to get all six equations because combining () and () will yield (), then () and () will yield (), then () and () will yield (), as well as would the combination of () and () as is explained in the following visual relation:
where the numbers represent the gas laws numbered above.
If you were to use the same method used above on 2 of the 3 laws on the vertices of one triangle that has a "O" inside it, you would get the third.
For example:
Change only pressure and volume first:
then only volume and temperature:
then as we can choose any value for , if we set , equation () becomes:
combining equations () and () yields , which is equation (), of which we had no prior knowledge until this derivation.
Theoretical
Kinetic theory
The ideal gas law can also be derived from first principles using the kinetic theory of gases, in which several simplifying assumptions are made, chief among which are that the molecules, or atoms, of the gas are point masses, possessing mass but no significant volume, and undergo only elastic collisions with each other and the sides of the container in which both linear momentum and kinetic energy are conserved.
First we show that the fundamental assumptions of the kinetic theory of gases imply that
Consider a container in the Cartesian coordinate system. For simplicity, we assume that a third of the molecules moves parallel to the -axis, a third moves parallel to the -axis and a third moves parallel to the -axis. If all molecules move with the same velocity , denote the corresponding pressure by . We choose an area on a wall of the container, perpendicular to the -axis. When time elapses, all molecules in the volume moving in the positive direction of the -axis will hit the area. There are molecules in a part of volume of the container, but only one sixth (i.e. a half of a third) of them moves in the positive direction of the -axis. Therefore, the number of molecules that will hit the area when the time elapses is .
When a molecule bounces off the wall of the container, it changes its momentum to . Hence the magnitude of change of the momentum of one molecule is . The magnitude of the change of momentum of all molecules that bounce off the area when time elapses is then . From and we get
We considered a situation where all molecules move with the same velocity . Now we consider a situation where they can move with different velocities, so we apply an "averaging transformation" to the above equation, effectively replacing by a new pressure and by the arithmetic mean of all squares of all velocities of the molecules, i.e. by Therefore
which gives the desired formula.
Using the Maxwell–Boltzmann distribution, the fraction of molecules that have a speed in the range to is , where
and denotes the Boltzmann constant. The root-mean-square speed can be calculated by
Using the integration formula
it follows that
from which we get the ideal gas law:
Statistical mechanics
Let q = (qx, qy, qz) and p = (px, py, pz) denote the position vector and momentum vector of a particle of an ideal gas, respectively. Let F denote the net force on that particle. Then (two times) the time-averaged kinetic energy of the particle is:
where the first equality is Newton's second law, and the second line uses Hamilton's equations and the equipartition theorem. Summing over a system of N particles yields
By Newton's third law and the ideal gas assumption, the net force of the system is the force applied by the walls of the container, and this force is given by the pressure P of the gas. Hence
where dS is the infinitesimal area element along the walls of the container. Since the divergence of the position vector q is
the divergence theorem implies that
where dV is an infinitesimal volume within the container and V is the total volume of the container.
Putting these equalities together yields
which immediately implies the ideal gas law for N particles:
where n = N/NA is the number of moles of gas and R = NAkB is the gas constant.
Other dimensions
For a d-dimensional system, the ideal gas pressure is:
where is the volume of the d-dimensional domain in which the gas exists. The dimensions of the pressure changes with dimensionality.
| Physical sciences | Thermodynamics | null |
59919 | https://en.wikipedia.org/wiki/Pheasant | Pheasant | Pheasants ( ) are birds of several genera within the family Phasianidae in the order Galliformes. Although they can be found all over the world in introduced (and captive) populations, the pheasant genera's native range is restricted to Eurasia. The classification "pheasant" is paraphyletic, as birds referred to as pheasants are included within both the subfamilies Phasianinae and Pavoninae, and in many cases are more closely related to smaller phasianids, grouse, and turkey (formerly classified in Perdicinae, Tetraoninae, and Meleagridinae) than to other pheasants.
Pheasants are characterised by strong sexual dimorphism, males being highly decorated with bright colours and adornments such as wattles. Males are usually larger than females and have longer tails. Males play no part in rearing the young.
A pheasant's call or cry can be recognised by the fact it sounds like a rusty sink or valve being turned.
Pheasants eat mostly seeds, grains, roots, and berries, while in the summer they take advantage of insects, fresh green shoots, spiders, earthworms, and snails. However, as an introduced species, in the UK they are a threat to endangered native adders.
The best-known is the common pheasant, which is widespread throughout the world, in introduced feral populations and in farm operations. Various other pheasant species are popular in aviaries, such as the golden pheasant (Chrysolophus pictus).
Etymology
According to the Oxford English Dictionary, the word "pheasant" ultimately comes from Phasis, the ancient name of what is now called the Rioni River in Georgia.
It passed from Greek to Latin to French (spelled with an initial "f") then to English, appearing for the first time in English around 1299.
Species in taxonomic order
This list is ordered to show presumed relationships between species.
Subfamily Phasianinae
Tribe Ithaginini
Blood pheasant (genus Ithaginis)
Blood pheasant (I. cruentus)
Tribe Pucrasiini
Koklass (genus Pucrasia)
Koklass pheasant (P. macrolopha)
Tribe Phasianini
Long-tailed pheasants (genus Syrmaticus)
Reeves's pheasant (S. reevesi)
Elliot's pheasant (S. ellioti)
Mrs. Hume's pheasant (S. humiae)
Mikado pheasant (S. mikado)
Copper pheasant (S. soemmerringi)
Ruffed pheasants (genus Chrysolophus)
Golden pheasant (C. pictus)
Lady Amherst's pheasant (C. amherstiae)
Typical pheasants (genus Phasianus)
Green pheasant (P. versicolor)
Common pheasant (P. colchicus)
Caucasus pheasants, Phasianus colchicus colchicus group
White-winged pheasants, Phasianus colchicus chrysomelas/principalis group
Prince of Wales pheasant, Phasianus colchicus principalis
Mongolian ring-necked pheasants or white-winged ring-necked pheasants, Phasianus colchicus mongolicus group
Tarim pheasants, Phasianus colchicus tarimensis group
Chinese ring-necked pheasants, Phasianus colchicus torquatus group
Taiwan pheasant, Phasianus colchicus formosanus
Cheer pheasant (genus Catreus)
Cheer pheasant (C. wallichi)
Gallopheasants (genus Lophura)
Kalij pheasant (L. leucomelanos)
White-crested kalij pheasant (L. l. hamiltoni)
Nepal kalij pheasant (L. l. leucomelanos)
Black-backed kalij pheasant (L. l. melanota)
Black kalij pheasant (L. l. moffitti)
Black-breasted kalij pheasant (L. l. lathami)
William's kalij pheasant (L. l. williamsi)
Oates' kalij pheasant (L. l. oatesi)
Crawfurd's kalij pheasant (L. l. crawfurdi)
Lineated kalij pheasant (L. l. lineata)
Silver pheasant (L. nycthemera)
Imperial pheasant (L. imperialis)
Edwards's pheasant (L. edwardsi)
Vietnamese pheasant (L. hatinhensis)
Swinhoe's pheasant (L. swinhoii)
Salvadori's pheasant (L. inornata)
Hoogerwerf's pheasant (L. i. hoogerwerfi)
Malayan crestless fireback (L. erythrophthalma)
Bornean crestless fireback (L. pyronota)
Bornean crested fireback (L. ignita)
Lesser Bornean crested fireback (L. i. ignita)
Greater Bornean crested fireback (L. i. nobilis)
Malayan crested fireback (L. rufa)
Siamese fireback (L. diardi)
Bulwer's pheasant (L. bulweri)
Eared pheasants (genus Crossoptilon)
White eared pheasant (C. crossoptilon)
Tibetan eared pheasant (C. harmani)
Brown eared pheasant (C. mantchuricum)
Blue eared pheasant (C. auritum)
Subfamily Pavoninae
Tribe Pavonini
Crested argus (genus Rheinardia)
Vietnamese crested argus (R. ocellata)
Malayan crested argus (R. nigrescens)
Great argus (genus Argusianus)
Great argus (A. argus)
Tribe Polyprectronini
Peacock-pheasants (genus Polyplectron)
Bronze-tailed peacock-pheasant (P. chalcurum)
Mountain peacock-pheasant (P. inopinatum)
Germain's peacock-pheasant (P. germaini)
Grey peacock-pheasant (P. bicalcaratum
Hainan peacock-pheasant (Polyplectron katsumatae)
Malayan peacock-pheasant (P. malacense)
Bornean peacock-pheasant (P. schleiermacheri)
Palawan peacock-pheasant (P. emphanum)
Previous classifications
Euplocamus and Gennceus are older names more or less corresponding to the current Lophura.
Euplocamus was used, for example, by Hume and Marshall in their Game Birds of India, Burmah and Ceylon (1879–1881).
Gennceus, was used, for example, by Frank Finn in Indian Sporting Birds (1915) and Game Birds of India and Asia (1911?).
These old genera were used for:
| Biology and health sciences | Galliformes | Animals |
59953 | https://en.wikipedia.org/wiki/Van%20Allen%20radiation%20belt | Van Allen radiation belt | The Van Allen radiation belt is a zone of energetic charged particles, most of which originate from the solar wind, that are captured by and held around a planet by that planet's magnetosphere. Earth has two such belts, and sometimes others may be temporarily created. The belts are named after James Van Allen, who published an article describing the belts in 1958.
Earth's two main belts extend from an altitude of about above the surface, in which region radiation levels vary. The belts are in the inner region of Earth's magnetic field. They trap energetic electrons and protons. Other nuclei, such as alpha particles, are less prevalent. Most of the particles that form the belts are thought to come from the solar wind while others arrive as cosmic rays. By trapping the solar wind, the magnetic field deflects those energetic particles and protects the atmosphere from destruction.
The belts endanger satellites, which must have their sensitive components protected with adequate shielding if they spend significant time near that zone. Apollo astronauts going through the Van Allen belts received a very low and harmless dose of radiation.
In 2013, the Van Allen Probes detected a transient, third radiation belt, which persisted for four weeks.
Discovery
Kristian Birkeland, Carl Størmer, Nicholas Christofilos, and Enrico Medi had investigated the possibility of trapped charged particles in 1895, forming a theoretical basis for the formation of radiation belts. The second Soviet satellite Sputnik 2 which had detectors designed by Sergei Vernov, followed by the US satellites Explorer 1 and Explorer 3, confirmed the existence of the belt in early 1958, later named after James Van Allen from the University of Iowa. The trapped radiation was first mapped by Explorer 4, Pioneer 3, and Luna 1.
The term Van Allen belts refers specifically to the radiation belts surrounding Earth; however, similar radiation belts have been discovered around other planets. The Sun does not support long-term radiation belts, as it lacks a stable, global dipole field. The Earth's atmosphere limits the belts' particles to regions above 200–1,000 km, (124–620 miles) while the belts do not extend past 8 Earth radii RE. The belts are confined to a volume which extends about 65° on either side of the celestial equator.
Research
The NASA Van Allen Probes mission aims at understanding (to the point of predictability) how populations of relativistic electrons and ions in space form or change in response to changes in solar activity and the solar wind.
NASA Institute for Advanced Concepts–funded studies have proposed magnetic scoops to collect antimatter that naturally occurs in the Van Allen belts of Earth, although only about 10 micrograms of antiprotons are estimated to exist in the entire belt.
The Van Allen Probes mission successfully launched on August 30, 2012. The primary mission was scheduled to last two years with expendables expected to last four. The probes were deactivated in 2019 after running out of fuel and are expected to deorbit during the 2030s. NASA's Goddard Space Flight Center manages the Living With a Star program—of which the Van Allen Probes were a project, along with Solar Dynamics Observatory (SDO). The Applied Physics Laboratory was responsible for the implementation and instrument management for the Van Allen Probes.
Radiation belts exist around other planets and moons in the solar system that have magnetic fields powerful and stable enough to sustain them. Radiation belts have been detected at Jupiter, Saturn, Uranus and Neptune through in-situ observations, such as by Galileo (spacecraft) and Juno (spacecraft) at Jupiter, Cassini–Huygens at Saturn and fly-bys from the Voyager program and Pioneer program. Observations of radio emissions from highly energetic particles that are trapped in a planets magnetic field have also been used to remotely detect radiation belts, including at Jupiter and at the ultracool dwarf LSR J1835+3259. It is possible that Mercury (planet) may be able to trap charged particles in its magnetic field, although its highly dynamic magnetosphere (which varies on the order of minutes ) may not be able to sustain stable radiation belts. Venus and Mars do not have radiation belts, as their magnetospheric configurations do not trap energetic charged particles in orbit around the planet.
Geomagnetic storms can cause electron density to increase or decrease relatively quickly (i.e., approximately one day or less). Longer-timescale processes determine the overall configuration of the belts. After electron injection increases electron density, electron density is often observed to decay exponentially. Those decay time constants are called "lifetimes." Measurements from the Van Allen Probe B's Magnetic Electron Ion Spectrometer (MagEIS) show long electron lifetimes (i.e., longer than 100 days) in the inner belt; short electron lifetimes of around one or two days are observed in the "slot" between the belts; and energy-dependent electron lifetimes of roughly five to 20 days are found in the outer belt.
Inner belt
The inner Van Allen Belt extends typically from an altitude of 0.2 to 2 Earth radii (L values of 1.2 to 3) or to above the Earth. In certain cases, when solar activity is stronger or in geographical areas such as the South Atlantic Anomaly, the inner boundary may decline to roughly 200 km above the Earth's surface. The inner belt contains high concentrations of electrons in the range of hundreds of keV and energetic protons with energies exceeding 100 MeV—trapped by the relatively strong magnetic fields in the region (as compared to the outer belt).
It is thought that proton energies exceeding 50 MeV in the lower belts at lower altitudes are the result of the beta decay of neutrons created by cosmic ray collisions with nuclei of the upper atmosphere. The source of lower energy protons is believed to be proton diffusion, due to changes in the magnetic field during geomagnetic storms.
Due to the slight offset of the belts from Earth's geometric center, the inner Van Allen belt makes its closest approach to the surface at the South Atlantic Anomaly.
In March 2014, a pattern resembling "zebra stripes" was observed in the radiation belts by the Radiation Belt Storm Probes Ion Composition Experiment (RBSPICE) onboard Van Allen Probes. The initial theory proposed in 2014 was that—due to the tilt in Earth's magnetic field axis—the planet's rotation generated an oscillating, weak electric field that permeates through the entire inner radiation belt. A 2016 study instead concluded that the zebra stripes were an imprint of ionospheric winds on radiation belts.
Outer belt
The outer belt consists mainly of high-energy (0.1–10 MeV) electrons trapped by the Earth's magnetosphere. It is more variable than the inner belt, as it is more easily influenced by solar activity. It is almost toroidal in shape, beginning at an altitude of 3 Earth radii and extending to 10 Earth radii (RE)— above the Earth's surface. Its greatest intensity is usually around 4 to 5 RE. The outer electron radiation belt is mostly produced by inward radial diffusion and local acceleration due to transfer of energy from whistler-mode plasma waves to radiation belt electrons. Radiation belt electrons are also constantly removed by collisions with Earth's atmosphere, losses to the magnetopause, and their outward radial diffusion. The gyroradii of energetic protons would be large enough to bring them into contact with the Earth's atmosphere. Within this belt, the electrons have a high flux and at the outer edge (close to the magnetopause), where geomagnetic field lines open into the geomagnetic "tail", the flux of energetic electrons can drop to the low interplanetary levels within about —a decrease by a factor of 1,000.
In 2014, it was discovered that the inner edge of the outer belt is characterized by a very sharp transition, below which highly relativistic electrons (> 5MeV) cannot penetrate. The reason for this shield-like behavior is not well understood.
The trapped particle population of the outer belt is varied, containing electrons and various ions. Most of the ions are in the form of energetic protons, but a certain percentage are alpha particles and O+ oxygen ions—similar to those in the ionosphere but much more energetic. This mixture of ions suggests that ring current particles probably originate from more than one source.
The outer belt is larger than the inner belt, and its particle population fluctuates widely. Energetic (radiation) particle fluxes can increase and decrease dramatically in response to geomagnetic storms, which are themselves triggered by magnetic field and plasma disturbances produced by the Sun. The increases are due to storm-related injections and acceleration of particles from the tail of the magnetosphere. Another cause of variability of the outer belt particle populations is the wave-particle interactions with various plasma waves in a broad range of frequencies.
On February 28, 2013, a third radiation belt—consisting of high-energy ultrarelativistic charged particles—was reported to be discovered. In a news conference by NASA's Van Allen Probe team, it was stated that this third belt is a product of coronal mass ejection from the Sun. It has been represented as a separate creation which splits the Outer Belt, like a knife, on its outer side, and exists separately as a storage container of particles for a month's time, before merging once again with the Outer Belt.
The unusual stability of this third, transient belt has been explained as due to a 'trapping' by the Earth's magnetic field of ultrarelativistic particles as they are lost from the second, traditional outer belt. While the outer zone, which forms and disappears over a day, is highly variable due to interactions with the atmosphere, the ultrarelativistic particles of the third belt are thought not to scatter into the atmosphere, as they are too energetic to interact with atmospheric waves at low latitudes. This absence of scattering and the trapping allows them to persist for a long time, finally only being destroyed by an unusual event, such as the shock wave from the Sun.
Flux values
In the belts, at a given point, the flux of particles of a given energy decreases sharply with energy.
At the magnetic equator, electrons of energies exceeding 5000 keV (resp. 5 MeV) have omnidirectional fluxes ranging from 1.2×106 (resp. 3.7×104) up to 9.4×109 (resp. 2×107) particles per square centimeter per second.
The proton belts contain protons with kinetic energies ranging from about 100 keV, which can penetrate 0.6 μm of lead, to over 400 MeV, which can penetrate 143 mm of lead.
Most published flux values for the inner and outer belts may not show the maximum probable flux densities that are possible in the belts. There is a reason for this discrepancy: the flux density and the location of the peak flux is variable, depending primarily on solar activity, and the number of spacecraft with instruments observing the belt in real time has been limited. The Earth has not yet experienced a solar storm of Carrington event intensity while spacecraft with the proper instruments have been available to observe the event.
Radiation levels in the belts would be dangerous to humans if they were exposed for an extended period of time. The Apollo missions minimised hazards for astronauts by sending spacecraft at high speeds through the thinner areas of the upper belts, bypassing inner belts completely, except for the Apollo 14 mission where the spacecraft traveled through the heart of the trapped radiation belts.
Antimatter confinement
In 2011, a study confirmed earlier speculation that the Van Allen belt could confine antiparticles. The Payload for Antimatter Matter Exploration and Light-nuclei Astrophysics (PAMELA) experiment detected levels of antiprotons orders of magnitude higher than are expected from normal particle decays while passing through the South Atlantic Anomaly. This suggests the Van Allen belts confine a significant flux of antiprotons produced by the interaction of the Earth's upper atmosphere with cosmic rays. The energy of the antiprotons has been measured in the range from 60 to 750 MeV.
The very high energy released in antimatter annihilation has led to proposals to harness these antiprotons for spacecraft propulsion. The concept relies on the development of antimatter collectors and containers.
Implications for space travel
Spacecraft travelling beyond low Earth orbit enter the zone of radiation of the Van Allen belts. Beyond the belts, they face additional hazards from cosmic rays and solar particle events. A region between the inner and outer Van Allen belts lies at 2 to 4 Earth radii and is sometimes referred to as the "safe zone".
Solar cells, integrated circuits, and sensors can be damaged by radiation. Geomagnetic storms occasionally damage electronic components on spacecraft. Miniaturization and digitization of electronics and logic circuits have made satellites more vulnerable to radiation, as the total electric charge in these circuits is now small enough so as to be comparable with the charge of incoming ions. Electronics on satellites must be hardened against radiation to operate reliably. The Hubble Space Telescope, among other satellites, often has its sensors turned off when passing through regions of intense radiation. A satellite shielded by 3 mm of aluminium in an elliptic orbit () passing the radiation belts will receive about 2,500 rem (25 Sv) per year. (For comparison, a full-body dose of 5 Sv is deadly.) Almost all radiation will be received while passing the inner belt.
The Apollo missions marked the first event where humans traveled through the Van Allen belts, which was one of several radiation hazards known by mission planners. The astronauts had low exposure in the Van Allen belts due to the short period of time spent flying through them.
Astronauts' overall exposure was actually dominated by solar particles once outside Earth's magnetic field. The total radiation received by the astronauts varied from mission-to-mission but was measured to be between 0.16 and 1.14 rads (1.6 and 11.4 mGy), much less than the standard of 5 rem (50 mSv) per year set by the United States Atomic Energy Commission for people who work with radioactivity.
Causes
It is generally understood that the inner and outer Van Allen belts result from different processes. The inner belt is mainly composed of energetic protons produced from the decay of so-called neutrons, which are themselves the result of cosmic ray collisions in the upper atmosphere. The outer Van Allen belt consists mainly of electrons. They are injected from the geomagnetic tail following geomagnetic storms, and are subsequently energized through wave-particle interactions.
In the inner belt, particles that originate from the Sun are trapped in the Earth's magnetic field. Particles spiral along the magnetic lines of flux as they move "latitudinally" along those lines. As particles move toward the poles, the magnetic field line density increases, and their "latitudinal" velocity is slowed and can be reversed, deflecting the particles back towards the equatorial region, causing them to bounce back and forth between the Earth's poles. In addition to both spiralling around and moving along the flux lines, the electrons drift slowly in an eastward direction, while the protons drift westward.
The gap between the inner and outer Van Allen belts is sometimes called the "safe zone" or "safe slot", and is the location of medium Earth orbits. The gap is caused by the VLF radio waves, which scatter particles in pitch angle, which adds new ions to the atmosphere. Solar outbursts can also dump particles into the gap, but those drain out in a matter of days. The VLF radio waves were previously thought to be generated by turbulence in the radiation belts, but recent work by J.L. Green of the Goddard Space Flight Center compared maps of lightning activity collected by the Microlab 1 spacecraft with data on radio waves in the radiation-belt gap from the IMAGE spacecraft; the results suggest that the radio waves are actually generated by lightning within Earth's atmosphere. The generated radio waves strike the ionosphere at the correct angle to pass through only at high latitudes, where the lower ends of the gap approach the upper atmosphere. These results are still being debated in the scientific community.
Proposed removal
Draining the charged particles from the Van Allen belts would open up new orbits for satellites and make travel safer for astronauts.
High Voltage Orbiting Long Tether, or HiVOLT, is a concept proposed by Russian physicist V. V. Danilov and further refined by Robert P. Hoyt and Robert L. Forward for draining and removing the radiation fields of the Van Allen radiation belts that surround the Earth.
Another proposal for draining the Van Allen belts involves beaming very-low-frequency (VLF) radio waves from the ground into the Van Allen belts.
Draining radiation belts around other planets has also been proposed, for example, before exploring Europa, which orbits within Jupiter's radiation belt.
As of 2024, it remains uncertain if there are any negative unintended consequences to removing these radiation belts.
| Physical sciences | Geophysics | Earth science |
59958 | https://en.wikipedia.org/wiki/Power%20series | Power series | In mathematics, a power series (in one variable) is an infinite series of the form
where represents the coefficient of the nth term and c is a constant called the center of the series. Power series are useful in mathematical analysis, where they arise as Taylor series of infinitely differentiable functions. In fact, Borel's theorem implies that every power series is the Taylor series of some smooth function.
In many situations, the center c is equal to zero, for instance for Maclaurin series. In such cases, the power series takes the simpler form
The partial sums of a power series are polynomials, the partial sums of the Taylor series of an analytic function are a sequence of converging polynomial approximations to the function at the center, and a converging power series can be seen as a kind of generalized polynomial with infinitely many terms. Conversely, every polynomial is a power series with only finitely many non-zero terms.
Beyond their role in mathematical analysis, power series also occur in combinatorics as generating functions (a kind of formal power series) and in electronic engineering (under the name of the Z-transform). The familiar decimal notation for real numbers can also be viewed as an example of a power series, with integer coefficients, but with the argument x fixed at . In number theory, the concept of p-adic numbers is also closely related to that of a power series.
Examples
Polynomial
Every polynomial of degree can be expressed as a power series around any center , where all terms of degree higher than have a coefficient of zero. For instance, the polynomial can be written as a power series around the center as
or around the center as
One can view power series as being like "polynomials of infinite degree", although power series are not polynomials in the strict sense.
Geometric series, exponential function and sine
The geometric series formula
which is valid for , is one of the most important examples of a power series, as are the exponential function formula
and the sine formula
valid for all real x. These power series are examples of Taylor series (or, more specifically, of Maclaurin series).
On the set of exponents
Negative powers are not permitted in an ordinary power series; for instance, is not considered a power series (although it is a Laurent series). Similarly, fractional powers such as are not permitted; fractional powers arise in Puiseux series. The coefficients must not depend on thus for instance is not a power series.
Radius of convergence
A power series is convergent for some values of the variable , which will always include since and the sum of the series is thus for . The series may diverge for other values of , possibly all of them. If is not the only point of convergence, then there is always a number with such that the series converges whenever and diverges whenever . The number is called the radius of convergence of the power series; in general it is given as
or, equivalently,
This is the Cauchy–Hadamard theorem; see limit superior and limit inferior for an explanation of the notation. The relation
is also satisfied, if this limit exists.
The set of the complex numbers such that is called the disc of convergence of the series. The series converges absolutely inside its disc of convergence and it converges uniformly on every compact subset of the disc of convergence.
For , there is no general statement on the convergence of the series. However, Abel's theorem states that if the series is convergent for some value such that , then the sum of the series for is the limit of the sum of the series for where is a real variable less than that tends to .
Operations on power series
Addition and subtraction
When two functions f and g are decomposed into power series around the same center c, the power series of the sum or difference of the functions can be obtained by termwise addition and subtraction. That is, if
and
then
The sum of two power series will have a radius of convergence of at least the smaller of the two radii of convergence of the two series, but possibly larger than either of the two. For instance it is not true that if two power series and have the same radius of convergence, then also has this radius of convergence: if and , for instance, then both series have the same radius of convergence of 1, but the series has a radius of convergence of 3.
Multiplication and division
With the same definitions for and , the power series of the product and quotient of the functions can be obtained as follows:
The sequence is known as the Cauchy product of the sequences and
For division, if one defines the sequence by
then
and one can solve recursively for the terms by comparing coefficients.
Solving the corresponding equations yields the formulae based on determinants of certain matrices of the coefficients of and
Differentiation and integration
Once a function is given as a power series as above, it is differentiable on the interior of the domain of convergence. It can be differentiated and integrated by treating every term separately since both differentiation and integration are linear transformations of functions:
Both of these series have the same radius of convergence as the original series.
Analytic functions
A function f defined on some open subset U of R or C is called analytic if it is locally given by a convergent power series. This means that every a ∈ U has an open neighborhood V ⊆ U, such that there exists a power series with center a that converges to f(x) for every x ∈ V.
Every power series with a positive radius of convergence is analytic on the interior of its region of convergence. All holomorphic functions are complex-analytic. Sums and products of analytic functions are analytic, as are quotients as long as the denominator is non-zero.
If a function is analytic, then it is infinitely differentiable, but in the real case the converse is not generally true. For an analytic function, the coefficients an can be computed as
where denotes the nth derivative of f at c, and . This means that every analytic function is locally represented by its Taylor series.
The global form of an analytic function is completely determined by its local behavior in the following sense: if f and g are two analytic functions defined on the same connected open set U, and if there exists an element such that for all , then for all .
If a power series with radius of convergence r is given, one can consider analytic continuations of the series, that is, analytic functions f which are defined on larger sets than and agree with the given power series on this set. The number r is maximal in the following sense: there always exists a complex number with such that no analytic continuation of the series can be defined at .
The power series expansion of the inverse function of an analytic function can be determined using the Lagrange inversion theorem.
Behavior near the boundary
The sum of a power series with a positive radius of convergence is an analytic function at every point in the interior of the disc of convergence. However, different behavior can occur at points on the boundary of that disc. For example:
Divergence while the sum extends to an analytic function: has radius of convergence equal to and diverges at every point of . Nevertheless, the sum in is , which is analytic at every point of the plane except for .
Convergent at some points divergent at others: has radius of convergence . It converges for , while it diverges for .
Absolute convergence at every point of the boundary: has radius of convergence , while it converges absolutely, and uniformly, at every point of due to Weierstrass M-test applied with the hyper-harmonic convergent series .
Convergent on the closure of the disc of convergence but not continuous sum: Sierpiński gave an example of a power series with radius of convergence , convergent at all points with , but the sum is an unbounded function and, in particular, discontinuous. A sufficient condition for one-sided continuity at a boundary point is given by Abel's theorem.
Formal power series
In abstract algebra, one attempts to capture the essence of power series without being restricted to the fields of real and complex numbers, and without the need to talk about convergence. This leads to the concept of formal power series, a concept of great utility in algebraic combinatorics.
Power series in several variables
An extension of the theory is necessary for the purposes of multivariable calculus. A power series is here defined to be an infinite series of the form
where is a vector of natural numbers, the coefficients are usually real or complex numbers, and the center and argument are usually real or complex vectors. The symbol is the product symbol, denoting multiplication. In the more convenient multi-index notation this can be written
where is the set of natural numbers, and so is the set of ordered n-tuples of natural numbers.
The theory of such series is trickier than for single-variable series, with more complicated regions of convergence. For instance, the power series is absolutely convergent in the set between two hyperbolas. (This is an example of a log-convex set, in the sense that the set of points , where lies in the above region, is a convex set. More generally, one can show that when c=0, the interior of the region of absolute convergence is always a log-convex set in this sense.) On the other hand, in the interior of this region of convergence one may differentiate and integrate under the series sign, just as one may with ordinary power series.
Order of a power series
Let be a multi-index for a power series . The order of the power series f is defined to be the least value such that there is aα ≠ 0 with , or if f ≡ 0. In particular, for a power series f(x) in a single variable x, the order of f is the smallest power of x with a nonzero coefficient. This definition readily extends to Laurent series.
| Mathematics | Sequences and series | null |
60023 | https://en.wikipedia.org/wiki/Airbag | Airbag | An airbag is a vehicle occupant-restraint system using a bag designed to inflate in milliseconds during a collision and then deflate afterwards. It consists of an airbag cushion, a flexible fabric bag, an inflation module, and an impact sensor. The purpose of the airbag is to provide a vehicle occupant with soft cushioning and restraint during a collision. It can reduce injuries between the flailing occupant and the vehicle's interior.
The airbag provides an energy-absorbing surface between the vehicle's occupants and a steering wheel, instrument panel, body pillar, headliner, and windshield. Modern vehicles may contain up to ten airbag modules in various configurations, including driver, passenger, side-curtain, seat-mounted, door-mounted, B and C-pillar mounted side-impact, knee bolster, inflatable seat belt, and pedestrian airbag modules.
During a crash, the vehicle's crash sensors provide crucial information to the airbag electronic controller unit (ECU), including collision type, angle, and severity of impact. Using this information, the airbag ECU's crash algorithm determines if the crash event meets the criteria for deployment and triggers various firing circuits to deploy one or more airbag modules within the vehicle. Airbag module deployments are activated through a pyrotechnic process designed to be used once as a supplemental restraint system for the vehicle's seat belt systems. Newer side-impact airbag modules consist of compressed-air cylinders that are triggered in the event of a side-on vehicle impact.
The first commercial designs were introduced in passenger automobiles during the 1970s, with limited success and caused some fatalities. Broad commercial adoption of airbags occurred in many markets during the late 1980s and early 1990s. Many modern vehicles now include six or more units.
Active vs. passive safety
Airbags are considered "passive" restraints and act as a supplement to "active" restraints. Because no action by a vehicle occupant is required to activate or use the airbag, it is considered a "passive" device. This is in contrast to seat belts, which are considered "active" devices because the vehicle occupant must act to enable them.
This terminology is not related to active and passive safety, which are, respectively, systems designed to prevent collisions in the first place, and systems designed to minimize the effects of collisions once they occur. In this use, a car Anti-lock braking system qualifies as an active-safety device, while both its seat belts and airbags qualify as passive-safety devices. Terminological confusion can arise from the fact that passive devices and systems—those requiring no input or action by the vehicle occupant—can operate independently in an active manner; an airbag is one such device. Vehicle safety professionals are generally careful in their use of language to avoid this sort of confusion. However, advertising principles sometimes prevent such semantic caution in the consumer marketing of safety features. Further confusing the terminology, the aviation safety community uses the terms "active" and "passive" in the opposite sense from the automotive industry.
History
Origins
The airbag "for the covering of aeroplane and other vehicle parts" traces its origins to a United States patent, submitted in 1919 by two dentists from Birmingham, Arthur Parrott and Harold Round. The patent was approved in 1920. Air-filled bladders were in use as early as 1951. The airbag specifically for automobile use is credited independently to the American John W. Hetrick, who filed for an airbag patent on 5 August 1952, that was granted #2,649,311 by the United States Patent Office on 18 August 1953. German engineer Walter Linderer, who filed German patent #896,312 on 6 October 1951, was issued on 12 November 1953, approximately three months after American John Hetrick. The airbags proposed by Hetrick and Linderer were based on compressed air released by a spring, bumper contact, or by the driver. Later research during the 1960s showed that compressed air could not inflate the mechanical airbags fast enough to ensure maximum safety, leading to the current chemical and electrical airbags. In patent applications, manufacturers sometimes use the term "inflatable occupant restraint systems".
Hetrick was an industrial engineer and member of the United States Navy. His airbag design, however, only came about when he combined his experiences working with navy torpedoes with his desire to protect his family on the road. Despite working with the major automobile manufacturers of his time, Hetrick was unable to attract investment. Although airbags are now required in every automobile sold in the United States, Hetrick's 1951 patent filing serves as an example of a "valuable" invention with little economic value to its inventor. Its first commercial use was not implemented until after the patent expired in 1971, at which point the airbag was installed in a few experimental Ford cars.
In 1964, a Japanese automobile engineer, Yasuzaburou Kobori (小堀保三郎), started developing an airbag "safety net" system. His design harnessed an explosive to inflate an airbag, for which he was later awarded patents in 14 countries. He died in 1975, before seeing the widespread adoption of airbag systems.
In 1967, a breakthrough in developing airbag crash sensors came when Allen K. Breed invented a ball-in-tube mechanism for crash detection. Under his system, an electromechanical sensor with a steel ball attached to a tube by a magnet would inflate an airbag in under 30 milliseconds. A small explosion of sodium azide was used instead of compressed air during inflation for the first time. Breed Corporation then marketed this innovation to Chrysler. A similar "Auto-Ceptor" crash-restraint, developed by the Eaton, Yale & Towne company for Ford, was soon also offered as an automatic safety system in the United States, while the Italian Eaton-Livia company offered a variant with localized air cushions.
In the early 1970s, General Motors began offering cars equipped with airbags, initially in government fleet-purchased 1973 Chevrolet Impala sedans. These cars came with a 1974-style Oldsmobile instrument panel and a unique steering wheel that contained the driver-side airbag. Two of these cars were crash tested after 20 years and the airbags deployed perfectly. An early example of the airbag cars survives as of 2009. GM's Oldsmobile Toronado was the first domestic U.S. vehicle to include a passenger airbag in 1973. General Motors marketed its first airbag modules under the "Air Cushion Restraint System" name, or ACRS. The automaker discontinued the option for its 1977 model year, citing a lack of consumer interest. Ford and GM then spent years lobbying against air-bag requirements, claiming that the devices were unfeasible and inappropriate. Chrysler made driver-side airbags standard on 1988 and 1989 models, but airbags did not become widespread in American cars until the early 1990s.
As a substitute for seat belts
Airbags for passenger cars were introduced in the United States in the 1970s. When seat-belt usage rates in the country were quite low compared to modern-day, Ford built experimental cars with airbags in 1971. Allstate operated a fleet of 200 Mercury Montereys and showed the reliability of airbags as well as their operation in crash testing, which also was promoted by the insurance company in popular magazine advertisements. General Motors followed in 1973 using full-sized Chevrolet vehicles. The early fleet of experimental GM vehicles equipped with airbags experienced seven fatalities, one of which was later suspected to have been caused by the airbag.
In 1974, GM made its ACRS system (which consisted of a padded lower dashboard and a passenger-side air bag) available as a regular production option (RPO code AR3) in full-sized Cadillac, Buick and Oldsmobile models. The GM cars from the 1970s equipped with ACRS had a driver-side airbag, and a driver-side knee restraint. The passenger-side airbag protected both front passengers, and unlike most modern systems, integrated a knee and torso cushion while also having a dual-stage deployment dictated by force of the impact. The cars equipped with ACRS had lap belts for all seating positions, but lacked shoulder belts. Shoulder belts were already mandatory in the United States on closed cars without airbags for the driver and outer front passenger, but GM chose to market its airbags as a substitute for shoulder belts. Prices for this option on Cadillac models were US$225 in 1974, $300 in 1975, and $340 in 1976 (US$ in dollars ).
The early development of airbags coincided with international interest in automobile safety legislation. Some safety experts advocated a performance-based occupant-protection standard rather than one mandating a particular technical solution (which could rapidly become outdated and prove to not be a cost-effective approach). Less emphasis was placed on other designs as countries successfully mandated seat belt restrictions, however.
As a supplemental restraint system
Frontal airbag
The auto industry and research and regulatory communities have moved away from their initial view of the airbag as a seat-belt replacement, and the bags are now nominally designated as supplemental restraint systems (SRS) or supplemental inflatable restraints.
In 1981, Mercedes-Benz introduced the airbag in West Germany as an option on its flagship saloon model, S-Class (W126). In the Mercedes system, the sensors automatically tensioned the seat belts to reduce occupants' motion on impact and then deployed the airbag on impact. This integrated the seat belts and the airbag into a restraint system, rather than the airbag being considered an alternative to the seat belt.
In 1987, the Porsche 944 Turbo became the first car to have driver and passenger airbags as standard equipment. The Porsche 944 and 944S had this as an available option. The same year also had the first airbag in a Japanese car, the Honda Legend.
In 1988, Chrysler became the first United States automaker to fit a driver-side airbag as standard equipment, which was offered in six different models. The following year, Chrysler became the first US auto manufacturer to offer driver-side airbags in all its new passenger models. Chrysler also began featuring the airbags in advertisements showing how the devices had saved lives that helped the public know the value of them and safety became a selling advantage in the late 1980s. All versions of the Chrysler minivans came with airbags starting for the 1991 model year. In 1993, The Lincoln Motor Company boasted that all vehicles in their model line were equipped with dual airbags, one for the driver's side and another for the passenger's side. The 1993 Jeep Grand Cherokee became the first SUV to offer a driver-side airbag when it was launched in 1992. Driver and passenger airbags became standard equipment in all Dodge Intrepid, Eagle Vision, and Chrysler Concorde sedans ahead of any safety regulations. Early 1993 saw the 4-millionth airbag-equipped Chrysler vehicle roll off the assembly line. In October 1993, the Dodge Ram became the first pickup truck with a standard driver-side airbag.
The first known collision between two airbag-equipped automobiles took place on 12 March 1990 in Virginia, USA. A 1989 Chrysler LeBaron crossed the center line and hit another 1989 Chrysler LeBaron in a head-on collision, causing both driver airbags to deploy. The drivers suffered only minor injuries despite extensive damage to the vehicles.
The United States Intermodal Surface Transportation Efficiency Act of 1991 required passenger cars and light trucks built after 1 September 1998 to have airbags for the driver and the front passenger. In the United States, NHTSA estimated that airbags had saved over 4,600 lives by 1 September 1999; however, the crash deployment experience of the early 1990s installations indicated that some fatalities and serious injuries were in fact caused by airbags. In 1998, NHTSA initiated new rules for advanced airbags that gave automakers more flexibility in devising effective technological solutions. The revised rules also required improved protection for occupants of different sizes regardless of whether they use seat belts, while minimizing the risk to infants, children, and other occupants caused by airbags.
In Europe, airbags were almost unheard of until the early 1990s. By 1991, four manufacturers – BMW, Honda, Mercedes-Benz, and Volvo – offered the airbag on some of their higher-end models, but shortly afterward, airbags became a common feature on more mainstream cars, with Ford and Vauxhall/Opel among the manufacturers to introduce the airbag to its model ranges in 1992. Citroën, Fiat, Nissan, Hyundai, Peugeot, Renault, and Volkswagen followed shortly afterwards.
By 1999, finding a new mass-market car without an airbag at least as optional equipment was difficult, and some late 1990s products, such as the Volkswagen Golf Mk4, also featured side airbags. The Peugeot 306 is one example of the European automotive mass-market evolution: starting in early 1993, most of these models did not even offer a driver's airbag as an option, but by 1999, even side airbags were available on several variants. Audi was late to offer airbag systems on a broader scale, since even in the 1994 model year, its popular models did not offer airbags. Instead, the German automaker until then relied solely on its proprietary cable-based procon-ten restraint system.
Variable force-deployment front airbags were developed to help minimize injury from the airbag itself.
The emergence of the airbag has contributed to a sharp decline in the number of deaths and serious injuries on the roads of Europe since 1990, and by 2010, the number of cars on European roads lacking an airbag represented a very small percentage of cars, mostly the remaining cars dating from the mid-1990s or earlier.
Many new cars in Latin America, including the Kia Rio, Kia Picanto, Hyundai Grand i10, Mazda 2, Chevrolet Spark and the Chevrolet Onix, are often sold without airbags, as neither airbags nor automatic braking systems in new cars are compulsory in many Latin American countries. Some require the installation of a minimum of only two airbags in new cars which many in this market have.
Shape of airbags
The Citroën C4 provided the first "shaped" driver airbag, made possible by this car's unusual fixed-hub steering wheel.
In 2019, Honda announced it would introduce a new front passenger airbag technology. Developed by Autoliv and Honda R&D in Ohio, United States, this new airbag design features three inflatable chambers connected across the front by a "noninflatable sail panel." The two outer chambers are larger than the middle chamber. When the airbag deploys, the sail panel cushions the occupant's head from the impact of hitting the airbag, and the three chambers hold the occupant's head in place, like a catcher's mitt. The goal of the tri-chamber airbag is to help "arrest high-speed movement" of the head, thereby reducing the likelihood of concussion injuries in a collision. The first vehicle to come with the tri-chamber airbag installed from the factory was in 2020 (for the 2021 model year) for the Acura TLX. Honda hopes that the new technology will soon make its way to all vehicles.
Rear airbag
Mercedes began offering rear passengers protection in frontal collisions in September 2020 (for the 2021 model year) for the Mercedes-Benz S-Class (W223). The W223 S-Class is the first car equipped with rear seat airbags that use gas to inflate supporting structures that unfold and extend a bag that fills with ambient air, instead of conventional fully gas-inflated airbags that are widely used in automotive airbag systems.
Side airbag
Essentially, two types of side airbags are commonly used today - the side-torso airbag and the side-curtain airbag. More recently, center airbags are becoming more common in the European market.
Most vehicles equipped with side-curtain airbags also include side-torso airbags. However, some, such as the Chevrolet Cobalt, 2007–09 model Chevrolet Silverado/GMC Sierra, and 2009–12 Dodge Ram do not feature the side-torso airbag.
From around 2000, side-impact airbags became commonplace on even low- to mid-range vehicles, such as the smaller-engined versions of the Ford Fiesta and Peugeot 206, and curtain airbags were also becoming regular features on mass-market cars. The Toyota Avensis, launched in 2003, was the first mass-market car to be sold in Europe with nine airbags.
Side torso airbag
Side-impact airbags or side-torso airbags are a category of airbags usually located in the seat or door panel, and inflate between the seat occupant and the door. These airbags are designed to reduce the risk of injury to the pelvic and lower abdomen regions. Most vehicles are now being equipped with different types of designs, to help reduce injury and ejection from the vehicle in rollover crashes. More recent side-airbag designs include a two-chamber system; a firmer lower chamber for the pelvic region and softer upper chamber for the ribcage.
Swedish company Autoliv AB was granted a patent on side-impact airbags, and they were first offered as an option in 1994 on the 1995 Volvo 850, and as standard equipment on all Volvo cars made after 1995.
In 1997, Saab introduced the first combined head and torso airbags with the launch of the Saab 9-5.
Some cars, such as the 2010 Volkswagen Polo Mk.5 have combined head- and torso-side airbags. These are fitted in the backrest of the front seats and protect the head and the torso.
Side tubular or curtain airbag
In 1997, the BMW 7 Series and 5 Series were fitted with tubular-shaped head side airbags (inflatable tubular structure), the "Head Protection System (HPS)" as standard equipment. This airbag was designed to offer head protection in side impact collisions and also maintained inflation for up to seven seconds for rollover protection. However, this tubular-shaped airbag design has been quickly replaced by an inflatable 'curtain' airbag.
In May 1998, Toyota began offering a side-curtain airbag deploying from the roof on the Progrés. In 1998, the Volvo S80 was given roof-mounted curtain airbags to protect both front and rear passengers. Curtain airbags were then made standard equipment on all new Volvo cars from 2000 except for the first-generation C70, which received an enlarged side-torso airbag that also protects the head of front-seat occupants. The second-generation C70 convertible received the world's first door-mounted, side-curtain airbags that deployed upwards.
Curtain airbags have been said to reduce brain injury or fatalities by up to 45% in a side impact with an SUV. These airbags come in various forms (e.g., tubular, curtain, door-mounted) depending on the needs of the application. Many recent SUVs and MPVs have a long inflatable curtain airbag that protects all rows of seats.
In many vehicles, the curtain airbags are programmed to deploy during some/all frontal impacts to manage passenger kinetics (e.g. head hitting B-pillar on the rebound), especially in offset crashes such as the IIHS's small overlap crash test.
Roll-sensing curtain airbag (RSCA)
Roll-sensing curtain airbags are designed to stay inflated for a longer duration of time, cover a larger proportion of the window, and be deployed in a roll-over crash. They offer protection to occupants' heads and help to prevent ejection. SUVs and pickups are more likely to be equipped with RSCAs due to their higher probability of rolling over and often a switch can disable the feature in case the driver wants to take the vehicle off-road.
Center airbag
In 2009, Toyota developed the first production rear-seat center airbag designed to reduce the severity of secondary injuries to rear passengers in a side collision. This system deploys from the rear center seat first appearing in on the Crown Majesta. In late 2012, General Motors with supplier Takata introduced a front center airbag; it deploys from the driver's seat.
Hyundai Motor Group announced its development of a center-side airbag on September 18, 2019, installed inside the driver's seat.
Some Volkswagen vehicles in 2022 equipped with center airbags include the ID.3 and the Golf.
The Polestar 2 also includes a center airbag.
With EuroNCAP updating its testing guidelines in 2020, European and Australian market vehicles increasingly use front-center airbags, rear torso airbags, and rear seat belt pre-tensioners.
Knee airbag
The second driver-side and separate knee airbag was used in the Kia Sportage SUV and has been standard equipment since then. The airbag is located beneath the steering wheel.
The Toyota Caldina introduced the first driver-side SRS knee airbag on the Japanese market in 2002. Toyota Avensis became the first vehicle sold in Europe equipped with a driver's knee airbag. The EuroNCAP reported on the 2003 Avensis, "There has been much effort to protect the driver's knees and legs and a knee airbag worked well." Since then certain models have also included front-passenger knee airbags, which deploy near or over the glove compartment in a crash. Knee airbags are designed to reduce leg injury. The knee airbag has become increasingly common since 2000.
Rear curtain airbag
In 2008, the new Toyota iQ microcar featured the first production rear-curtain shield airbag to protect the rear occupants' heads in the event of a rear-end impact.
Seat cushion airbag
Another feature of the Toyota iQ was a seat-cushion airbag in the passenger seat to prevent the pelvis from diving below the lap belt during a frontal impact or submarining. Later Toyota models such as the Yaris added the feature to the driver's seat, as well.
Seat-belt airbag
The seat-belt airbag is designed to better distribute the forces experienced by a buckled person in a crash using an increased seat belt area. This is done to reduce possible injuries to the rib cage or chest of the belt wearer.
2010: Ford Explorer and 2013 Ford Flex: optional rear seat belt airbags; standard on the 2013 Lincoln MKT
2010: Lexus LFA had seat belt airbags for driver and passenger
2013: Mercedes-Benz S-Class (W222) has rear seat belt bags
2014: Ford Mondeo Mk IV has optional rear seat belt airbags for the two outer seats
Cessna Aircraft also introduced seat belt airbags. They are as of 2003 standard on the 172, 182, and 206.
Pedestrian airbag
Airbag(s) mounted to the exterior of vehicles, so-called "pedestrian airbags", are designed to reduce injuries in the event of a vehicle to a pedestrian collision. When a collision is detected the airbag will deploy and cover hard areas, such as a-pillars and bonnet edges, before they can be struck by the pedestrian.
When introduced in 2012 the Volvo V40 included the world's first pedestrian airbag as standard. As a result, the V40 ranked highest (88%) in the EuroNCAP's pedestrian tests.
Manufacturers
Suppliers of SRS airbags include Autoliv, Daicel, TRW, and JSS (which owns Breed, Key Safety Systems, and Takata). The majority of impact sensors of airbags are manufactured by the Lanka Harness Company.
Operation
The airbags in the vehicle are controlled by a central airbag control unit (ACU), a specific type of ECU. The ACU monitors several related sensors within the vehicle, including accelerometers, impact sensors, side (door) pressure sensors, wheel speed sensors, gyroscopes, brake pressure sensors, and seat occupancy sensors. Oftentimes, ACUs log this—and other—sensor data in a circular buffer and record it to onboard non-volatile memory, to provide a snapshot of the crash event for investigators. As such, an ACU frequently functions as the vehicle's event data recorder; not all EDRs are ACUs, and not all ACUs include EDR features. An ACU typically includes capacitors within its circuitry, so that the module remains powered and able to deploy the airbags if the vehicle's battery connection to the ACU is severed during a crash.
The bag itself and its inflation mechanism are concealed within the steering wheel boss (for the driver), or the dashboard (for the front passenger), behind plastic flaps or doors that are designed to tear open under the force of the bag inflating. Once the requisite threshold has been reached or exceeded, the airbag control unit will trigger the ignition of a gas generator propellant to rapidly inflate a fabric bag. As the vehicle occupant collides with and squeezes the bag, the gas escapes in a controlled manner through small vent holes. The airbag's volume and the size of the vents in the bag are tailored to each vehicle type, to spread out the deceleration of (and thus force experienced by) the occupant over time and over the occupant's body, compared to a seat belt alone.
The signals from the various sensors are fed into the airbag control unit, which determines from them the angle of impact, the severity, or the force of the crash, along with other variables. Depending on the result of these calculations, the ACU may also deploy various additional restraint devices, such as seat belt pre-tensioners, and/or airbags (including frontal bags for driver and front passenger, along with seat-mounted side bags, and "curtain" airbags which cover the side glass). Each restraint device is typically activated with one or more pyrotechnic devices, commonly called an initiator or electric match. The electric match, which consists of an electrical conductor wrapped in a combustible material, activates with a current pulse between 1 and 3 amperes in less than 2 milliseconds. When the conductor becomes hot enough, it ignites the combustible material, which initiates the gas generator. In a seat belt pre-tensioner, this hot gas is used to drive a piston that pulls the slack out of the seat belt. In an airbag, the initiator is used to ignite solid propellant inside the airbag inflator. The burning propellant generates inert gas which rapidly inflates the airbag in approximately 20 to 30 milliseconds. An airbag must inflate quickly to be fully inflated by the time the forward-traveling occupant reaches its outer surface. Typically, the decision to deploy an airbag in a frontal crash is made within 15 to 30 milliseconds after the onset of the crash, and both the driver and passenger airbags are fully inflated within approximately 60–80 milliseconds after the first moment of vehicle contact. If an airbag deploys too late or too slowly, the risk of occupant injury from contact with the inflating airbag may increase. Since more distance typically exists between the passenger and the instrument panel, the passenger airbag is larger and requires more gas to fill it.
Older airbag systems contained a mixture of sodium azide (NaN3), KNO3, and SiO2. A typical driver-side airbag contains approximately 50–80 g of NaN3, with the larger passenger-side airbag containing about 250 g. Within about 40 milliseconds of impact, all these components react in three separate reactions that produce nitrogen gas. The reactions, in order, are as follows.
2 NaN3 → 2 Na + 3 N2 (g)
10 Na + 2 KNO3 → K2O + 5 Na2O + N2 (g)
K2O + Na2O + 2 SiO2 → K2SiO3 + Na2SiO3
The first two reactions create 4 molar equivalents of nitrogen gas, and the third converts the remaining reactants to relatively inert potassium silicate and sodium silicate. The reason that KNO3 is used rather than something like NaNO3 is because it is less hygroscopic. The materials used in this reaction must not be hygroscopic because absorbed moisture can de-sensitize the system and cause the reaction to fail.
The particle size of the initial reactants is important to reliable operation. The NaN3 and KNO3 must be between 10 and 20 µm, while the SiO2 must be between 5 and 10 µm.
There are ongoing efforts to find alternative compounds so that airbags have less toxic reactants. The reaction of the Sr complex nitrate, (Sr(NH2NHCONHNH2)∙(NO3)2) of carbohydrazide (SrCDH) with various oxidizing agents resultS in the evolution of N2 and CO2 gases. Using KBrO3 as the oxidizing agent resulted in the most vigorous reaction as well as the lowest initial temperature of the reaction. The N2 and CO2 gases evolved made up 99% of all gases evolved. Nearly all the starting materials will not decompose until reaching temperatures of 500 °C or higher, so this could be a viable option as an airbag gas generator.
In a patent containing another plausible alternative to NaN3 driven airbags, the gas-generating materials involved the use of guanidine nitrate, 5-aminotetrazole, bitetrazole dihydrate, nitroimidazole, and basic copper nitrate. It was found that these non-azide reagents allowed for a less toxic, lower combustion temperature reaction, and more easily disposable airbag inflation system.
Front airbags normally do not protect the occupants during side, rear, or rollover collisions. Since airbags deploy only once and deflate quickly after the initial impact, they will not be beneficial during a subsequent collision. Safety belts help reduce the risk of injury in many types of crashes. They help to properly position occupants to maximize the airbag's benefits and they help restrain occupants during the initial and any following collisions.
In vehicles equipped with a rollover sensing system, accelerometers, and gyroscopes are used to sense the onset of a rollover event. If a rollover event is determined to be imminent, side-curtain airbags are deployed to help protect the occupant from contact with the side of the vehicle interior, and also to help prevent occupant ejection as the vehicle rolls over.
Triggering conditions
Airbags are designed to deploy in frontal and near-frontal collisions more severe than a threshold defined by the regulations governing vehicle construction in whatever particular market the vehicle is intended for: United States regulations require deployment in crashes at least equivalent in deceleration to a barrier collision, or similarly, striking a parked car of similar size across the full front of each vehicle at about twice the speed. International regulations are performance-based, rather than technology-based, so airbag deployment threshold is a function of overall vehicle design.
Unlike crash tests into barriers, real-world crashes typically occur at angles other than directly into the front of the vehicle, and the crash forces usually are not evenly distributed across the front of the vehicle. Consequently, the relative speed between a striking and struck vehicle required to deploy the airbag in a real-world crash can be much higher than an equivalent barrier crash. Because airbag sensors measure deceleration, the vehicle speed is not a good indicator of whether an airbag should be deployed. Airbags can deploy due to the vehicle's undercarriage striking a low object protruding above the roadway due to the resulting deceleration.
The airbag sensor is a MEMS accelerometer, which is a small integrated circuit with integrated micromechanical elements. The microscopic mechanical element moves in response to rapid deceleration, and this motion causes a change in capacitance, which is detected by the electronics on the chip that then sends a signal to fire the airbag. The most common MEMS accelerometer in use is the ADXL-50 by Analog Devices, but there are other MEMS manufacturers as well.
Initial attempts using mercury switches did not work well. Before MEMS, the primary system used to deploy airbags was called a "rolamite". A rolamite is a mechanical device, consisting of a roller suspended within a tensioned band. As a result of the particular geometry and material properties used, the roller is free to translate with little friction or hysteresis. This device was developed at Sandia National Laboratories. Rolamite and similar macro-mechanical devices were used in airbags until the mid-1990s after which they were universally replaced with MEMS.
Nearly all airbags are designed to automatically deploy in the event of a vehicle fire when temperatures reach . This safety feature, often termed auto-ignition, helps to ensure that such temperatures do not cause an explosion of the entire airbag module.
Today, airbag triggering algorithms are much more complex, being able to adapt the deployment speed to the crash conditions, and prevent unnecessary deployments. The algorithms are considered valuable intellectual property. Experimental algorithms may take into account such factors as the weight of the occupant, the seat location, and seat belt use, as well as even attempt to determine if a baby seat is present.
Inflation
When the frontal airbags are to deploy, a signal is sent to the inflator unit within the airbag control unit. An igniter starts a rapid chemical reaction generating primarily nitrogen gas (N2) to fill the airbag making it deploy through the module cover. Some airbag technologies use compressed nitrogen or argon gas with a pyrotechnic operated valve ("hybrid gas generator"), while other technologies use various energetic propellants. Although propellants containing the highly toxic sodium azide (NaN3) were common in early inflator designs, little to no toxic sodium azide has been found on used airbags.
The azide-containing pyrotechnic gas generators contain a substantial amount of the propellant. The driver-side airbag would contain a canister containing about 50 grams of sodium azide. The passenger side container holds about 200 grams of sodium azide.
The alternative propellants may incorporate, for example, a combination of nitroguanidine, phase-stabilized ammonium nitrate (NH4NO3) or another nonmetallic oxidizer, and a nitrogen-rich fuel different from azide (e.g. tetrazoles, triazoles, and their salts). The burn rate modifiers in the mixture may be an alkaline metal nitrate (NO3-) or nitrite (NO2-), dicyanamide or its salts, sodium borohydride (NaBH4), etc. The coolants and slag formers may be e.g. clay, silica, alumina, glass, etc. Other alternatives are e.g. nitrocellulose based propellants (which have high gas yield but bad storage stability, and their oxygen balance requires secondary oxidation of the reaction products to avoid buildup of carbon monoxide), or high-oxygen nitrogen-free organic compounds with inorganic oxidizers (e.g., di or tricarboxylic acids with chlorates (ClO3-) or perchlorates (ClO4-) and eventually metallic oxides; the nitrogen-free formulation avoids formation of toxic nitrogen oxides).
From the onset of the crash, the entire deployment and inflation process is about 0.04 seconds. Because vehicles change speed so quickly in a crash, airbags must inflate rapidly to reduce the risk of the occupant hitting the vehicle's interior.
Variable-force deployment
Advanced airbag technologies are being developed to tailor airbag deployment to the severity of the crash, the size, and posture of the vehicle occupant, belt usage, and how close that person is to the actual airbag. Many of these systems use multi-stage inflators that deploy less forcefully in stages in moderate crashes than in very severe crashes. Occupant sensing devices let the airbag control unit know if someone is occupying a seat adjacent to an airbag, the mass/weight of the person, whether a seat belt or child restraint is being used, and whether the person is forward in the seat and close to the airbag. Based on this information and crash severity information, the airbag is deployed at either a high force level, a less forceful level, or not at all.
Adaptive airbag systems may utilize multi-stage airbags to adjust the pressure within the airbag. The greater the pressure within the airbag, the more force the airbag will exert on the occupants as they come in contact with it. These adjustments allow the system to deploy the airbag with a moderate force for most collisions; reserving the maximum force airbag only for the severest of collisions. Additional sensors to determine the location, weight or relative size of the occupants may also be used. Information regarding the occupants and the severity of the crash are used by the airbag control unit, to determine whether airbags should be suppressed or deployed, and if so, at various output levels.
Post-deployment
A chemical reaction produces a burst of nitrogen to inflate the bag. Once an airbag deploys, deflation begins immediately as the gas escapes through vent(s) in the fabric (or, as it is sometimes called, the cushion) and cools. Deployment is frequently accompanied by the release of dust-like particles, and gases in the vehicle's interior (called effluent). Most of this dust consists of cornstarch, french chalk, or talcum powder, which are used to lubricate the airbag during deployment.
Newer designs produce effluent primarily consisting of harmless talcum powder/cornstarch and nitrogen gas. In older designs using an azide-based propellant (usually NaN3), varying amounts of sodium hydroxide nearly always are initially present. In small amounts this chemical can cause minor irritation to the eyes and/or open wounds; however, with exposure to air, it quickly turns into sodium bicarbonate (baking soda). However, this transformation is not 100% complete, and invariably leaves residual amounts of hydroxide ions from NaOH. Depending on the type of airbag system, potassium chloride may also be present.
For most people, the only effect the dust may produce is some minor irritation of the throat and eyes. Generally, minor irritations only occur when the occupant remains in the vehicle for many minutes with the windows closed and no ventilation. However, some people with asthma may develop a potentially lethal asthmatic attack from inhaling the dust.
Because of the airbag exit flap design of the steering wheel boss and dashboard panel, these items are not designed to be recoverable if an airbag deploys, meaning that they have to be replaced if the vehicle has not been written off in a collision. Moreover, the dust-like particles and gases can cause irreparable cosmetic damage to the dashboard and upholstery, meaning that minor collisions that result in the deployment of airbags can be costly, even if there are no injuries and there is only minor damage to the vehicle structure.
Regulatory specifications
United States
On 11 July 1984, the United States government amended Federal Motor Vehicle Safety Standard 208 (FMVSS 208) to require cars produced after 1 April 1989 to be equipped with a passive restraint for the driver. An airbag or an automatic seat belt would meet the requirements of the standard. Airbag introduction was stimulated by the National Highway Traffic Safety Administration. However, airbags were not mandatory on light trucks until 1997.
In 1998, FMVSS 208 was amended to require dual front airbags, and reduced-power, second-generation airbags were also mandated. This was due to the injuries caused by first-generation airbags, though FMVSS 208 continues to require that bags be engineered and calibrated to be able to "save" the life of an unbelted 50th-percentile size and weight "male" crash test dummy. The technical performance and validation requirements for the inflator assembly used in airbag modules are specified in SAE USCAR 24–2.
Europe
Some countries outside North America adhere to internationalized European ECE vehicle and equipment regulations rather than the United States Federal Motor Vehicle Safety Standards. ECE airbags are generally smaller and inflate less forcefully than United States airbags because the ECE specifications are based on belted crash test dummies.
The Euro NCAP vehicle safety rating encourages manufacturers to take a comprehensive approach to occupant safety; a good rating can only be achieved by combining airbags with other safety features. Almost every new car sold in Europe is equipped with front and side airbags, but in the European Union in 2020 and in the United Kingdom, and most other developed countries there is no direct legal requirement for new cars to feature airbags.
Central and South America
Ecuador requires dual front airbags in new car models since 2013.
Since January 2014, except for micro vehicles, all new cars made or imported in Argentina must have front airbags.
Since 1 January 2014, all new cars sold in Brazil must have dual front airbags.
Since July 2014, all new cars sold in Uruguay must have dual front airbags.
Since December 2016, all new cars sold in Chile must have dual front airbags.
Since 1 January 2017, all cars made or imported in Colombia must have dual front airbags.
Since 1 January 2020, all new cars sold in Mexico must have dual front airbags.
India
On 5 March 2021, the Indian Ministry of Road Transport and Highways mandated that all new vehicle models introduced in India after 1 April 2021 have dual front airbags; the regulation also requires that all existing models be equipped with dual front airbags by 31 August 2021.
India also mandated that all passenger vehicles sold after October 2023 must have a minimum of six airbags.
Maintenance
Inadvertent airbag deployment while the vehicle is being serviced can result in severe injury, and an improperly installed or defective airbag unit may not operate or perform as intended. Those servicing a vehicle, as well as first responders, must exercise extreme awareness, as many airbag control systems may remain powered for roughly 30 minutes after disconnecting the vehicle's battery. Some countries impose restrictions on the sale, transport, handling, and service of airbags and system components. In Germany, airbags are regulated as harmful explosives; only mechanics with special training are allowed to service airbag systems.
Some automakers (such as Mercedes-Benz) call for the replacement of undeployed airbags after a certain time to ensure their reliability in a collision. One example is the 1992 S500, which has an expiry date sticker attached to the door pillar. Some Škoda vehicles indicate an expiry date of 14 years from the date of manufacture. In this case, replacement would be uneconomic as the car would have negligible value at 14 years old, far less than the cost of fitting new airbags. Volvo has stated that "airbags do not require replacement during the lifetime of the vehicle," though this cannot be taken as a guarantee on the device.
Limitations
Although the millions of installed airbags in use have an excellent safety record, some limitations on their ability to protect car occupants exist. The original implementation of front airbags did little to protect against side collisions, which can be more dangerous than frontal collisions because the protective crumple zone in front of the passenger compartment is completely bypassed. Side airbags and protective airbag curtains are increasingly being required in modern vehicles to protect against this very common category of collisions.
Airbags are designed to deploy once only, so are ineffective if any further collisions occur after an initial impact. Multiple impacts may occur during rollovers or other incidents involving multiple collisions, such as many multivehicle collisions.
An extremely dangerous situation occurs during "underride collisions", in which a passenger vehicle collides with the rear of a tractor-trailer lacking a rear underride guard, or hits the side of such a trailer not equipped with a side underride guard. The platform bed of a typical trailer is roughly at the head height of a seated adult occupant of a typical passenger car. This means not much of a barrier exists between a head and the edge of the trailer platform, except a glass windshield. In an underride collision, the car's crush zones designed to absorb collision energy are completely bypassed, and the airbags may not deploy in time because the car does not decelerate appreciably until the windshield and roof pillars have already impacted the trailer bed. Even delayed inflation of airbags may be useless because of major intrusion into the passenger space, leaving occupants at high risk of major head trauma or decapitation in even low-speed collisions. Western European standards for underride guards have been stricter than North American standards, which typically have allowed grandfathering of older equipment that may still be on the road for decades.
Typical airbag systems are completely disabled by turning off the ignition key. Unexpected turnoffs usually also disable the engine, power steering, and power brakes, and can be the direct cause of a collision. If a violent collision occurs, the disabled airbags will not deploy to protect vehicle occupants. In 2014, General Motors admitted to concealing information about fatal collisions caused by defective ignition switches that would abruptly shut down a car (including its airbags). Between 13 and 74 deaths have been directly attributed to this defect, depending on how the fatalities are classified.
Injuries and fatalities
Under some rare conditions, airbags can injure and in some very rare instances kill vehicle occupants. To provide crash protection for occupants not wearing seat belts, United States airbag designs trigger much more forcefully than airbags designed to the international ECE standards used in most other countries. Recent "smart" airbag controllers can recognize if a seat belt is used, and alter the airbag cushion deployment parameters accordingly.
In 1990, the first automotive fatality attributed to an airbag was reported. TRW produced the first gas-inflated airbag in 1994, with sensors and low inflation-force bags becoming common soon afterward. Dual-depth (also known as dual-stage) airbags appeared on passenger cars in 1998. By 2005, deaths related to airbags had declined, with no adult deaths and two child deaths attributed to airbags that year. However, injuries remain fairly common in collisions with airbag deployment.
Serious injuries are less common, but severe or fatal injuries can occur to vehicle occupants very near an airbag or in direct contact when it deploys. Such injuries may be sustained by unconscious drivers slumped over the steering wheel, unrestrained or improperly restrained occupants who slide forward in the seat during precrash braking, and properly belted drivers sitting very close to the steering wheel. A good reason for the driver not to cross hands over the steering wheel, a rule taught to most learner drivers, but quickly forgotten by most, is that an airbag deployment while negotiating a turn may result in the driver's hand(s) being driven forcefully into his or her face, exacerbating any injuries from the airbag alone.
Improvements in sensing and gas-generator technology have allowed the development of third-generation airbag systems that can adjust their deployment parameters to the size, weight, position, and restraint status of the occupant. These improvements have demonstrated a reduced injury risk factor for small adults and children, who had an increased risk of injury with first-generation airbag systems.
One model of airbags made by the Takata Corporation used ammonium nitrate–based gas-generating compositions in airbag inflators instead of the more stable, but more expensive compound tetrazole. The ammonium nitrate-based inflators have a flaw where old inflators with long-term exposure to hot and humid climate conditions could rupture during deployment, projecting metal shards through the airbag and into the driver. As of December 2022, the defect has caused 33 deaths worldwide, with up to 24 in the U.S. and the remaining in Australia and Malaysia. The National Highway Traffic Safety Administration (NHTSA) recalled over 33 million vehicles in May 2015, and fined Takata $70 million in November 2015. Toyota, Mazda, and Honda have said that they will not use ammonium-nitrate inflators. In June 2017, Takata filed for bankruptcy.
Airbag fatality statistics
From 1990 to 2000, the United States NHTSA identified 175 fatalities caused by airbags. Most of these (104) have been children. About 3.3 million airbag deployments have occurred during that interval, and the agency estimates more than 6,377 lives were saved and countless injuries were prevented.
A rear-facing infant restraint put in the front seat of a vehicle places an infant's head close to the airbag, which can cause severe head injuries or death if the airbag deploys. Some modern cars include a switch to disable the front-passenger airbag, in case a child-supporting seat is used there (although not in Australia, where rear-facing child seats are prohibited in the front where an airbag is fitted).
In vehicles with side airbags, it is dangerous for occupants to lean against the windows, doors, and pillars, or to place objects between themselves and the side of the vehicle. Articles hung from a vehicle's clothes hanger hooks can be hazardous if the vehicle's side-curtain airbags are deployed. A seat-mounted airbag may also cause internal injury if the occupant leans against the door.
Aerospace and military applications
The aerospace industry and the United States government have applied airbag technologies for many years. NASA and United States Department of Defense have incorporated airbag systems in various aircraft and spacecraft applications as early as the 1960s.
Spacecraft airbag landing systems
The first use of airbags for landing were Luna 9 and Luna 13. As with later missions, these would use the airbags to bounce along the surface, absorbing landing energy. The Mars Pathfinder lander employed an innovative airbag landing system, supplemented with aerobraking, parachute, and solid rocket landing thrusters. This prototype successfully tested the concept, and the two Mars Exploration Rover Mission landers employed similar landing systems. The Beagle 2 Mars lander also tried to use airbags for landing; the landing was successful, and the lander touched down safely, but several of the spacecraft's solar panels failed to deploy, thereby disabling the spacecraft.
The Boeing Starliner uses six airbags to cushion the ground landing of the crewed capsule.
Aircraft airbag landing systems
Airbags have also been used on military fixed-wing aircraft, such as the escape crew capsule of the F-111 Aardvark.
Occupant protection
The United States Army has incorporated airbags in its UH-60A/L Black Hawk and OH-58D Kiowa Warrior helicopter fleets. The Cockpit Air Bag System (CABS) consists of forward and lateral airbags, and an inflatable tubular structure (on the OH-58D only) with an Electronic Crash Sensor Unit (ECSU). The CABS system was developed by the United States Army Aviation Applied Technology Directorate, through a contract with Simula Safety Systems (now BAE Systems). It is the first conventional airbag system for occupant injury prevention (worldwide) designed and developed and placed in service for an aircraft, and the first specifically for helicopter applications.
Other uses
In the mid-1970s, the UK Transport Research Laboratory tested several types of motorcycle airbags. In 2006 Honda introduced the first production motorcycle airbag safety system on its Gold Wing motorcycle. Honda claims that sensors in the front forks can detect a severe frontal collision and decide when to deploy the airbag, absorbing some of the forward energy of the rider and reducing the velocity at which the rider may be thrown from the motorcycle.
More commonly, air bag vests – either integrated into the motorcyclist's jacket or worn over it – have started to become used by regular riders on the street. MotoGP has made it compulsory since 2018 for riders to wear suits with integrated airbags.
Similarly, companies such as Helite and Hit-Air have commercialized equestrian airbags, which attach to the saddle and are worn by the rider. Other sports, particularly skiing and snowboarding, have started introducing airbag safety mechanisms.
| Technology | Basics_7 | null |
60088 | https://en.wikipedia.org/wiki/Roentgenium | Roentgenium | Roentgenium () is a synthetic chemical element; it has symbol Rg and atomic number 111. It is extremely radioactive and can only be created in a laboratory. The most stable known isotope, roentgenium-282, has a half-life of 130 seconds, although the unconfirmed roentgenium-286 may have a longer half-life of about 10.7 minutes. Roentgenium was first created in December 1994 by the GSI Helmholtz Centre for Heavy Ion Research near Darmstadt, Germany. It is named after the physicist Wilhelm Röntgen (also spelled Roentgen), who discovered X-rays. Only a few roentgenium atoms have ever been synthesized, and they have no practical application.
In the periodic table, it is a d-block transactinide element. It is a member of the 7th period and is placed in the group 11 elements, although no chemical experiments have been carried out to confirm that it behaves as the heavier homologue to gold in group 11 as the ninth member of the 6d series of transition metals. Roentgenium is calculated to have similar properties to its lighter homologues, copper, silver, and gold, although it may show some differences from them.
Introduction
History
Official discovery
Roentgenium was first synthesized by an international team led by Sigurd Hofmann at the Gesellschaft für Schwerionenforschung (GSI) in Darmstadt, Germany, on December 8, 1994. The team bombarded a target of bismuth-209 with accelerated nuclei of nickel-64 and detected three nuclei of the isotope roentgenium-272:
+ → +
This reaction had previously been conducted at the Joint Institute for Nuclear Research in Dubna (then in the Soviet Union) in 1986, but no atoms of 272Rg had then been observed. In 2001, the IUPAC/IUPAP Joint Working Party (JWP) concluded that there was insufficient evidence for the discovery at that time. The GSI team repeated their experiment in 2002 and detected three more atoms. In their 2003 report, the JWP decided that the GSI team should be acknowledged for the discovery of this element.
Naming
Using Mendeleev's nomenclature for unnamed and undiscovered elements, roentgenium should be known as eka-gold. In 1979, IUPAC published recommendations according to which the element was to be called unununium (with the corresponding symbol of Uuu), a systematic element name as a placeholder, until the element was discovered (and the discovery then confirmed) and a permanent name was decided on. Although widely used in the chemical community on all levels, from chemistry classrooms to advanced textbooks, the recommendations were mostly ignored among scientists in the field, who called it element 111, with the symbol of E111, (111) or even simply 111.
The name roentgenium (Rg) was suggested by the GSI team in 2004, to honor the German physicist Wilhelm Conrad Röntgen, the discoverer of X-rays. This name was accepted by IUPAC on November 1, 2004.
Isotopes
Roentgenium has no stable or naturally occurring isotopes. Several radioactive isotopes have been synthesized in the laboratory, either by fusion of the nuclei of lighter elements or as intermediate decay products of heavier elements. Nine different isotopes of roentgenium have been reported with atomic masses 272, 274, 278–283, and 286 (283 and 286 unconfirmed), two of which, roentgenium-272 and roentgenium-274, have known but unconfirmed metastable states. All of these decay through alpha decay or spontaneous fission, though 280Rg may also have an electron capture branch.
Stability and half-lives
All roentgenium isotopes are extremely unstable and radioactive; in general, the heavier isotopes are more stable than the lighter. The most stable known roentgenium isotope, 282Rg, is also the heaviest known roentgenium isotope; it has a half-life of 100 seconds. The unconfirmed 286Rg is even heavier and appears to have an even longer half-life of about 10.7 minutes, which would make it one of the longest-lived superheavy nuclides known; likewise, the unconfirmed 283Rg appears to have a long half-life of about 5.1 minutes. The isotopes 280Rg and 281Rg have also been reported to have half-lives over a second. The remaining isotopes have half-lives in the millisecond range.
The missing isotopes between 274Rg and 278Rg are too light to be produced by hot fusion and too heavy to be produced by cold fusion. A possible synthesis method is to populate them from above, as daughters of nihonium or moscovium isotopes that can be produced by hot fusion. The isotopes 283Rg and 284Rg could be synthesised using charged-particle evaporation, using the 238U+48Ca reaction where a proton is evaporated alongside some neutrons.
Predicted properties
Other than nuclear properties, no properties of roentgenium or its compounds have been measured; this is due to its extremely limited and expensive production and the fact that roentgenium (and its parents) decays very quickly. Properties of roentgenium metal remain unknown and only predictions are available.
Chemical
Roentgenium is the ninth member of the 6d series of transition metals. Calculations on its ionization potentials and atomic and ionic radii are similar to that of its lighter homologue gold, thus implying that roentgenium's basic properties will resemble those of the other group 11 elements, copper, silver, and gold; however, it is also predicted to show several differences from its lighter homologues.
Roentgenium is predicted to be a noble metal. The standard electrode potential of 1.9 V for the Rg3+/Rg couple is greater than that of 1.5 V for the Au3+/Au couple. Roentgenium's predicted first ionisation energy of 1020 kJ/mol almost matches that of the noble gas radon at 1037 kJ/mol. Its predicted second ionization energy, 2070 kJ/mol, is almost the same as that of silver. Based on the most stable oxidation states of the lighter group 11 elements, roentgenium is predicted to show stable +5 and +3 oxidation states, with a less stable +1 state. The +3 state is predicted to be the most stable. Roentgenium(III) is expected to be of comparable reactivity to gold(III), but should be more stable and form a larger variety of compounds. Gold also forms a somewhat stable −1 state due to relativistic effects, and it has been suggested roentgenium may do so as well: nevertheless, the electron affinity of roentgenium is expected to be around , significantly lower than gold's value of , so roentgenides may not be stable or even possible.
The 6d orbitals are destabilized by relativistic effects and spin–orbit interactions near the end of the fourth transition metal series, thus making the high oxidation state roentgenium(V) more stable than its lighter homologue gold(V) (known only in gold pentafluoride, Au2F10) as the 6d electrons participate in bonding to a greater extent. The spin-orbit interactions stabilize molecular roentgenium compounds with more bonding 6d electrons; for example, is expected to be more stable than , which is expected to be more stable than . The stability of is homologous to that of ; the silver analogue is unknown and is expected to be only marginally stable to decomposition to and F2. Moreover, Rg2F10 is expected to be stable to decomposition, exactly analogous to the Au2F10, whereas Ag2F10 should be unstable to decomposition to Ag2F6 and F2. Gold heptafluoride, AuF7, is known as a gold(V) difluorine complex AuF5·F2, which is lower in energy than a true gold(VII) heptafluoride would be; RgF7 is instead calculated to be more stable as a true roentgenium(VII) heptafluoride, although it would be somewhat unstable, its decomposition to Rg2F10 and F2 releasing a small amount of energy at room temperature. Roentgenium(I) is expected to be difficult to obtain. Gold readily forms the cyanide complex , which is used in its extraction from ore through the process of gold cyanidation; roentgenium is expected to follow suit and form .
The probable chemistry of roentgenium has received more interest than that of the two previous elements, meitnerium and darmstadtium, as the valence s-subshells of the group 11 elements are expected to be relativistically contracted most strongly at roentgenium. Calculations on the molecular compound RgH show that relativistic effects double the strength of the roentgenium–hydrogen bond, even though spin–orbit interactions also weaken it by . The compounds AuX and RgX, where X = F, Cl, Br, O, Au, or Rg, were also studied. Rg+ is predicted to be the softest metal ion, even softer than Au+, although there is disagreement on whether it would behave as an acid or a base. In aqueous solution, Rg+ would form the aqua ion [Rg(H2O)2]+, with an Rg–O bond distance of 207.1 pm. It is also expected to form Rg(I) complexes with ammonia, phosphine, and hydrogen sulfide.
Physical and atomic
Roentgenium is expected to be a solid under normal conditions and to crystallize in the body-centered cubic structure, unlike its lighter congeners which crystallize in the face-centered cubic structure, due to its being expected to have different electron charge densities from them. It should be a very heavy metal with a density of around 22–24 g/cm3; in comparison, the densest known element that has had its density measured, osmium, has a density of 22.61 g/cm3. The atomic radius of roentgenium is expected to be around 138 pm.
Experimental chemistry
Unambiguous determination of the chemical characteristics of roentgenium has yet to have been established due to the low yields of reactions that produce roentgenium isotopes. For chemical studies to be carried out on a transactinide, at least four atoms must be produced, the half-life of the isotope used must be at least 1 second, and the rate of production must be at least one atom per week. Even though the half-life of 282Rg, the most stable confirmed roentgenium isotope, is 100 seconds, long enough to perform chemical studies, another obstacle is the need to increase the rate of production of roentgenium isotopes and allow experiments to carry on for weeks or months so that statistically significant results can be obtained. Separation and detection must be carried out continuously to separate out the roentgenium isotopes and allow automated systems to experiment on the gas-phase and solution chemistry of roentgenium, as the yields for heavier elements are predicted to be smaller than those for lighter elements. However, the experimental chemistry of roentgenium has not received as much attention as that of the heavier elements from copernicium to livermorium, despite early interest in theoretical predictions due to relativistic effects on the ns subshell in group 11 reaching a maximum at roentgenium. The isotopes 280Rg and 281Rg are promising for chemical experimentation and may be produced as the granddaughters of the moscovium isotopes 288Mc and 289Mc respectively; their parents are the nihonium isotopes 284Nh and 285Nh, which have already received preliminary chemical investigations.
| Physical sciences | Group 11 | Chemistry |
60162 | https://en.wikipedia.org/wiki/Tidal%20locking | Tidal locking | Tidal locking between a pair of co-orbiting astronomical bodies occurs when one of the objects reaches a state where there is no longer any net change in its rotation rate over the course of a complete orbit. In the case where a tidally locked body possesses synchronous rotation, the object takes just as long to rotate around its own axis as it does to revolve around its partner. For example, the same side of the Moon always faces Earth, although there is some variability because the Moon's orbit is not perfectly circular. Usually, only the satellite is tidally locked to the larger body. However, if both the difference in mass between the two bodies and the distance between them are relatively small, each may be tidally locked to the other; this is the case for Pluto and Charon, and for Eris and Dysnomia. Alternative names for the tidal locking process are gravitational locking, captured rotation, and spin–orbit locking.
The effect arises between two bodies when their gravitational interaction slows a body's rotation until it becomes tidally locked. Over many millions of years, the interaction forces changes to their orbits and rotation rates as a result of energy exchange and heat dissipation. When one of the bodies reaches a state where there is no longer any net change in its rotation rate over the course of a complete orbit, it is said to be tidally locked. The object tends to stay in this state because leaving it would require adding energy back into the system. The object's orbit may migrate over time so as to undo the tidal lock, for example, if a giant planet perturbs the object.
There is ambiguity in the use of the terms 'tidally locked' and 'tidal locking', in that some scientific sources use it to refer exclusively to 1:1 synchronous rotation (e.g. the Moon), while others include non-synchronous orbital resonances in which there is no further transfer of angular momentum over the course of one orbit (e.g. Mercury). In Mercury's case, the planet completes three rotations for every two revolutions around the Sun, a 3:2 spin–orbit resonance. In the special case where an orbit is nearly circular and the body's rotation axis is not significantly tilted, such as the Moon, tidal locking results in the same hemisphere of the revolving object constantly facing its partner.
Regardless of which definition of tidal locking is used, the hemisphere that is visible changes slightly due to variations in the locked body's orbital velocity and the inclination of its rotation axis over time.
Mechanism
Consider a pair of co-orbiting objects, A and B. The change in rotation rate necessary to tidally lock body B to the larger body A is caused by the torque applied by A's gravity on bulges it has induced on B by tidal forces.
The gravitational force from object A upon B will vary with distance, being greatest at the nearest surface to A and least at the most distant. This creates a gravitational gradient across object B that will distort its equilibrium shape slightly. The body of object B will become elongated along the axis oriented toward A, and conversely, slightly reduced in dimension in directions orthogonal to this axis. The elongated distortions are known as tidal bulges. (For the solid Earth, these bulges can reach displacements of up to around .) When B is not yet tidally locked, the bulges travel over its surface due to orbital motions, with one of the two "high" tidal bulges traveling close to the point where body A is overhead. For large astronomical bodies that are nearly spherical due to self-gravitation, the tidal distortion produces a slightly prolate spheroid, i.e. an axially symmetric ellipsoid that is elongated along its major axis. Smaller bodies also experience distortion, but this distortion is less regular.
The material of B exerts resistance to this periodic reshaping caused by the tidal force. In effect, some time is required to reshape B to the gravitational equilibrium shape, by which time the forming bulges have already been carried some distance away from the A–B axis by B's rotation. Seen from a vantage point in space, the points of maximum bulge extension are displaced from the axis oriented toward A. If B's rotation period is shorter than its orbital period, the bulges are carried forward of the axis oriented toward A in the direction of rotation, whereas if B's rotation period is longer, the bulges instead lag behind.
Because the bulges are now displaced from the A–B axis, A's gravitational pull on the mass in them exerts a torque on B. The torque on the A-facing bulge acts to bring B's rotation in line with its orbital period, whereas the "back" bulge, which faces away from A, acts in the opposite sense. However, the bulge on the A-facing side is closer to A than the back bulge by a distance of approximately B's diameter, and so experiences a slightly stronger gravitational force and torque. The net resulting torque from both bulges, then, is always in the direction that acts to synchronize B's rotation with its orbital period, leading eventually to tidal locking.
Orbital changes
The angular momentum of the whole A–B system is conserved in this process, so that when B slows down and loses rotational angular momentum, its orbital angular momentum is boosted by a similar amount (there are also some smaller effects on A's rotation). This results in a raising of B's orbit about A in tandem with its rotational slowdown. For the other case where B starts off rotating too slowly, tidal locking both speeds up its rotation, and lowers its orbit.
Locking of the larger body
The tidal locking effect is also experienced by the larger body A, but at a slower rate because B's gravitational effect is weaker due to B's smaller mass. For example, Earth's rotation is gradually being slowed by the Moon, by an amount that becomes noticeable over geological time as revealed in the fossil record. Current estimations are that this (together with the tidal influence of the Sun) has helped lengthen the Earth day from about 6 hours to the current 24 hours (over about 4.5 billion years). Currently, atomic clocks show that Earth's day lengthens, on average, by about 2.3 milliseconds per century. Given enough time, this would create a mutual tidal locking between Earth and the Moon. The length of Earth's day would increase and the length of a lunar month would also increase. Earth's sidereal day would eventually have the same length as the Moon's orbital period, about 47 times the length of the Earth day at present. However, Earth is not expected to become tidally locked to the Moon before the Sun becomes a red giant and engulfs Earth and the Moon.
For bodies of similar size the effect may be of comparable size for both, and both may become tidally locked to each other on a much shorter timescale. An example is the dwarf planet Pluto and its satellite Charon. They have already reached a state where Charon is visible from only one hemisphere of Pluto and vice versa.
Eccentric orbits
For orbits that do not have an eccentricity close to zero, the rotation rate tends to become locked with the orbital speed when the body is at periapsis, which is the point of strongest tidal interaction between the two objects. If the orbiting object has a companion, this third body can cause the rotation rate of the parent object to vary in an oscillatory manner. This interaction can also drive an increase in orbital eccentricity of the orbiting object around the primary – an effect known as eccentricity pumping.
In some cases where the orbit is eccentric and the tidal effect is relatively weak, the smaller body may end up in a so-called spin–orbit resonance, rather than being tidally locked. Here, the ratio of the rotation period of a body to its own orbital period is some simple fraction different from 1:1. A well known case is the rotation of Mercury, which is locked to its own orbit around the Sun in a 3:2 resonance. This results in the rotation speed roughly matching the orbital speed around perihelion.
Many exoplanets (especially the close-in ones) are expected to be in spin–orbit resonances higher than 1:1. A Mercury-like terrestrial planet can, for example, become captured in a 3:2, 2:1, or 5:2 spin–orbit resonance, with the probability of each being dependent on the orbital eccentricity.
Occurrence
Moons
All twenty known moons in the Solar System that are large enough to be round are tidally locked with their primaries, because they orbit very closely and tidal force increases rapidly (as a cubic function) with decreasing distance. On the other hand, most of the irregular outer satellites of the giant planets (e.g. Phoebe), which orbit much farther away than the large well-known moons, are not tidally locked.
Pluto and Charon are an extreme example of a tidal lock. Charon is a relatively large moon in comparison to its primary and also has a very close orbit. This results in Pluto and Charon being mutually tidally locked. Pluto's other moons are not tidally locked; Styx, Nix, Kerberos, and Hydra all rotate chaotically due to the influence of Charon. Similarly, and Dysnomia are mutually tidally locked. and Vanth might also be mutually tidally locked, but the data is not conclusive.
The tidal locking situation for asteroid moons is largely unknown, but closely orbiting binaries are expected to be tidally locked, as well as contact binaries.
Earth's Moon
Earth's Moon's rotation and orbital periods are tidally locked with each other, so no matter when the Moon is observed from Earth, the same hemisphere of the Moon is always seen. Most of the far side of the Moon was not seen until 1959, when photographs of most of the far side were transmitted from the Soviet spacecraft Luna 3.
When Earth is observed from the Moon, Earth does not appear to move across the sky. It remains in the same place while showing nearly all its surface as it rotates on its axis.
Despite the Moon's rotational and orbital periods being exactly locked, about 59 percent of the Moon's total surface may be seen with repeated observations from Earth, due to the phenomena of libration and parallax. Librations are primarily caused by the Moon's varying orbital speed due to the eccentricity of its orbit: this allows up to about 6° more along its perimeter to be seen from Earth. Parallax is a geometric effect: at the surface of Earth observers are offset from the line through the centers of Earth and Moon; this accounts for about a 1° difference in the Moon's surface which can be seen around the sides of the Moon when comparing observations made during moonrise and moonset.
Planets
It was thought for some time that Mercury was in synchronous rotation with the Sun. This was because whenever Mercury was best placed for observation, the same side faced inward. Radar observations in 1965 demonstrated instead that Mercury has a 3:2 spin–orbit resonance, rotating three times for every two revolutions around the Sun, which results in the same positioning at those observation points. Modeling has demonstrated that Mercury was captured into the 3:2 spin–orbit state very early in its history, probably within 10–20 million years after its formation.
The 583.92-day interval between successive close approaches of Venus to Earth is equal to 5.001444 Venusian solar days, making approximately the same face visible from Earth at each close approach. Whether this relationship arose by chance or is the result of some kind of tidal locking with Earth is unknown.
The exoplanet Proxima Centauri b discovered in 2016 which orbits around Proxima Centauri, is almost certainly tidally locked, expressing either synchronized rotation or a 3:2 spin–orbit resonance like that of Mercury.
One form of hypothetical tidally locked exoplanets are eyeball planets, which in turn are divided into "hot" and "cold" eyeball planets.
Stars
Close binary stars throughout the universe are expected to be tidally locked with each other, and extrasolar planets that have been found to orbit their primaries extremely closely are also thought to be tidally locked to them. An unusual example, confirmed by MOST, may be Tau Boötis, a star that is probably tidally locked by its planet Tau Boötis b. If so, the tidal locking is almost certainly mutual.
Timescale
An estimate of the time for a body to become tidally locked can be obtained using the following formula:
where
is the initial spin rate expressed in radians per second,
is the semi-major axis of the motion of the satellite around the planet (given by the average of the periapsis and apoapsis distances),
is the moment of inertia of the satellite, where is the mass of the satellite and is the mean radius of the satellite,
is the dissipation function of the satellite,
is the gravitational constant,
is the mass of the planet (i.e., the object being orbited), and
is the tidal Love number of the satellite.
and are generally very poorly known except for the Moon, which has . For a really rough estimate it is common to take (perhaps conservatively, giving overestimated locking times), and
where
is the density of the satellite
is the surface gravity of the satellite
is the rigidity of the satellite. This can be roughly taken as 3 N/m2 for rocky objects and 4 N/m2 for icy ones.
Even knowing the size and density of the satellite leaves many parameters that must be estimated (especially ω, Q, and μ), so that any calculated locking times obtained are expected to be inaccurate, even to factors of ten. Further, during the tidal locking phase the semi-major axis may have been significantly different from that observed nowadays due to subsequent tidal acceleration, and the locking time is extremely sensitive to this value.
Because the uncertainty is so high, the above formulas can be simplified to give a somewhat less cumbersome one. By assuming that the satellite is spherical, , and it is sensible to guess one revolution every 12 hours in the initial non-locked state (most asteroids have rotational periods between about 2 hours and about 2 days)
with masses in kilograms, distances in meters, and in newtons per meter squared; can be roughly taken as 3 N/m2 for rocky objects and 4 N/m2 for icy ones.
There is an extremely strong dependence on semi-major axis .
For the locking of a primary body to its satellite as in the case of Pluto, the satellite and primary body parameters can be swapped.
One conclusion is that, other things being equal (such as and ), a large moon will lock faster than a smaller moon at the same orbital distance from the planet because grows as the cube of the satellite radius . A possible example of this is in the Saturn system, where Hyperion is not tidally locked, whereas the larger Iapetus, which orbits at a greater distance, is. However, this is not clear cut because Hyperion also experiences strong driving from the nearby Titan, which forces its rotation to be chaotic.
The above formulae for the timescale of locking may be off by orders of magnitude, because they ignore the frequency dependence of . More importantly, they may be inapplicable to viscous binaries (double stars, or double asteroids that are rubble), because the spin–orbit dynamics of such bodies is defined mainly by their viscosity, not rigidity.
List of known tidally locked bodies
Solar System
All the bodies below are tidally locked, and all but Mercury are moreover in synchronous rotation. (Mercury is tidally locked, but not in synchronous rotation.)
Extra-solar
The most successful detection methods of exoplanets (transits and radial velocities) suffer from a clear observational bias favoring the detection of planets near the star; thus, 85% of the exoplanets detected are inside the tidal locking zone, which makes it difficult to estimate the true incidence of this phenomenon. Tau Boötis is known to be locked to the close-orbiting giant planet Tau Boötis b.
Bodies likely to be locked
Solar System
Based on comparison between the likely time needed to lock a body to its primary, and the time it has been in its present orbit (comparable with the age of the Solar System for most planetary moons), a number of moons are thought to be locked. However their rotations are not known or not known enough. These are:
Probably locked to Saturn
Daphnis
Aegaeon
Methone
Anthe
Pallene
Helene
Polydeuces
Probably locked to Uranus
Cordelia
Ophelia
Bianca
Cressida
Desdemona
Juliet
Portia
Rosalind
Cupid
Belinda
Perdita
Puck
Mab
Probably locked to Neptune
Naiad
Thalassa
Despina
Galatea
Larissa
Probably mutually tidally locked
Orcus and Vanth
Extrasolar
Gliese 581c, Gliese 581g, Gliese 581b, and Gliese 581e may be tidally locked to their parent star Gliese 581. Gliese 581d is almost certainly captured either into the 2:1 or the 3:2 spin–orbit resonance with the same star.
All planets in the TRAPPIST-1 system are likely to be tidally locked.
| Physical sciences | Celestial mechanics | null |
60209 | https://en.wikipedia.org/wiki/Ascaris%20lumbricoides | Ascaris lumbricoides | Ascaris lumbricoides is a large parasitic roundworm of the genus Ascaris. It is the most common parasitic worm in humans. An estimated 807 million–1.2 billion people are infected with A. lumbricoides worldwide. People living in tropical and subtropical countries are at greater risk of infection. Infection by Ascaris lumbricoides is known as ascariasis.
It has been proposed that Ascaris lumbricoides and Ascaris suum (pig roundworm) are the same species.
Life cycle
Ascaris lumbricoides, a roundworm, infects humans via the fecal-oral route. Eggs released by adult females are shed in feces. Unfertilized eggs are often observed in fecal samples but never become infective. Fertilized eggs embryonate and become infectious after 18 days to several weeks in soil, depending on the environmental conditions (optimum: moist, warm, shaded soil).
Infection occurs when a human swallows water or food contaminated with embryonated eggs. In the duodenum, a single rhabditiform larva hatches from each of the ingested eggs. The larvae then penetrate the mucosa and submucosa and enter the venules or lymphatic vessels. From there, the larvae then pass through the heart to enter the pulmonary circulation. The larvae then break through the walls of the pulmonary capillaries to enter the alveoli. The juvenile worms then migrate from the alveoli, through the bronchioles and bronchi, and into the trachea. An acute inflammatory reaction can occur if some of the worms get lost during this migration process and accumulate in other organs of the body.
Once in the trachea, the worms are coughed up into the pharynx and then swallowed again, after which they pass through the stomach and into the small intestine, where they mature into adult worms. The adult worms begin producing fertilized eggs within 60–65 days of being swallowed; females produce as many as 200,000 eggs per day for 12–18 months. These fertilized eggs become infectious after two weeks in soil; they can persist in soil for 10 years or more.
It might seem odd that the worms end up in the same place where they began. One hypothesis to account for this behavior is that the migration mimics an intermediate host, which would be required for juveniles of an ancestral form to develop to the third stage. Another possibility is that tissue migration enables faster growth and larger size, which increases reproductive capacity. The eggs have a lipid layer which makes them resistant to the effects of acids and alkalis, as well as other chemicals.
Morphology
Ascaris lumbricoides is characterized by its great size. Males are in diameter and long. The male's posterior end is curved ventrally and has a bluntly pointed tail. Females are wide and long. The vulva is located in the anterior end and accounts for about one-third of its body length. Uteri may contain up to 27 million eggs at a time, with 200,000 being laid per day. Fertilized eggs are oval to round in shape and are long and wide with a thick outer shell. Unfertilized eggs measure long and wide.
Epidemiology
An estimated 807 million–1.2 billion people are infected with A. lumbricoides worldwide. While infection occurs throughout most of the world, ascariasis is most common in sub-Saharan Africa, the Americas, China, and east Asia. Although the prevalence is low in the United States, ascariasis is still endemic in the southeastern United States due to the temperature and humid climate.
A. lumbricoides eggs are extremely resistant to strong chemicals, desiccation, and low temperatures. The eggs can remain viable in soil for months or even years. Eggs of A. lumbricoides have been identified in coprolites in the Americas, Europe, Africa, the Middle East, and New Zealand, the oldest ones being more than 24,000 years old.
Infections
Infections with these parasites are more common where sanitation is poor, and raw human feces are used as fertilizer.
Symptoms
Often, no symptoms are presented with a minor A. lumbricoides infection, the inevitable consequence being the e.g. once a year passage of such clearly visible worm(s) on close inspection. In the case of bad infections symptoms commonly include bloody sputum, cough, fever, abdominal discomfort, intestinal ulcer(s), as well as a less commonly missed passing of the quite long worms. Ascariasis is the most common cause of Löffler's syndrome worldwide. Accompanying pathological symptoms include pulmonary infiltration, eosinophilia (symptoms of the overabundance of eosinophils in the blood such as asthma and allergic reactions), and a diagnostic symptom is, aside from standard microscopy of stools, radiographic opacities. One study has observed increases in fertility in infected women, in a similar vein to good diet and exercise, but with all of the pathological negatives and discomforts the disease carries with it, varying from host to host and again with diet.
Distribution
Ascaris lumbricoides is primarily distributed in tropical and subtropical regions around the world, particularly in areas with poor sanitation and hygiene practices. It is most prevalent in sub-Saharan Africa, Southeast Asia (including countries like India, Bangladesh, and Indonesia), and parts of Latin America, where inadequate sanitation infrastructure and the use of human faeces as fertilizer contribute to its spread.
Prevention
Preventing any fecal-borne disease requires educated hygienic habits/culture and effective fecal treatment systems. This is particularly important with A. lumbricoides because its eggs are one of the most difficult pathogens to kill (second only to prions), and the eggs commonly survive 1–3 years. A. lumbricoides lives in the intestine where it lays eggs. Infection occurs when the eggs, too small to be seen by the unaided eye, are eaten. The eggs may get onto vegetables when improperly processed human feces of infected people are used as fertilizer for food crops. Infection may occur when food is handled without removing or killing the eggs on the hands, clothes, hair, raw vegetables/fruit, or cooked food that is (re)infected by handlers, containers, etc. Bleach does not readily kill A. lumbricoides eggs, but it will remove their sticky film, to allow the eggs to be rinsed away. A. lumbricoides eggs can be reduced by hot composting methods, but to completely kill them may require rubbing alcohol, iodine, specialized chemicals, cooking heat, or "unusually" hot composting (for example, over for 24 hours).
Treatment
Control of roundworm infections is based on treatment with medication, improved sanitation and health education. This usually takes around three days.
History
Giant intestinal roundworms have been known since antiquity. In 1758 Linnaeus named them Ascaris lumbricoides. For many centuries, they were thought to arise by spontaneous generation. In 1855, Ascaris eggs were found in human faeces by Henry Ransom in England then this was described in the literature two years later by Casimir-Joseph Davaine in France. Attempts to infect animals by feeding them eggs were unsuccessful. In 1886, Salvatore Calandruccio in Italy successfully infected a boy to whom he had given 150 eggs. Battista Grassi published this information without giving any acknowledgement to Calandruccio. Development was thought to occur directly within the bowel lumen but in Francis Stewart in Hong Kong in 1916 fed eggs to rats, then later mice, and found infective larvae in the faeces and in the lungs but no mature worms. In 1918, Sadao Yoshida ingested larvae recovered from the trachea of a guinea pig, then found eggs in his own stools 76 days later. In 1922, Shimesu Koino ingested 2,000 Ascaris lumbricoides eggs, found larvae in his sputum a few days later, then after 50 days took an anthelmintic and recovered 667 immature Ascaris lumbricoides, thus confirming the life cycle.
| Biology and health sciences | Ecdysozoa | Animals |
60244 | https://en.wikipedia.org/wiki/Dugong | Dugong | The dugong (; Dugong dugon) is a marine mammal. It is one of four living species of the order Sirenia, which also includes three species of manatees. It is the only living representative of the once-diverse family Dugongidae; its closest modern relative, Steller's sea cow (Hydrodamalis gigas), was hunted to extinction in the 18th century.
The dugong is the only sirenian in its range, which spans the waters of some 40 countries and territories throughout the Indo-West Pacific. The dugong is largely dependent on seagrass communities for subsistence and is thus restricted to the coastal habitats that support seagrass meadows, with the largest dugong concentrations typically occurring in wide, shallow, protected areas such as bays, mangrove channels, the waters of large inshore islands and inter-reefal waters. The northern waters of Australia between Shark Bay and Moreton Bay are believed to be the dugong's contemporary stronghold.
Like all modern sirenians, the dugong has a fusiform body with no dorsal fin or hind limbs. The forelimbs or flippers are paddle-like. The dugong is easily distinguishable from the manatees by its fluked, dolphin-like tail; moreover, it possesses a unique skull and teeth. Its snout is sharply downturned, an adaptation for feeding in benthic seagrass communities. The molar teeth are simple and peg-like, unlike the more elaborate molar dentition of manatees.
The dugong has been hunted for thousands of years for its meat and oil. Traditional hunting still has great cultural significance in several countries in its modern range, particularly northern Australia and the Pacific Islands. The dugong's current distribution is fragmented, and many populations are believed to be close to extinction. The IUCN lists the dugong as a species vulnerable to extinction, while the Convention on International Trade in Endangered Species limits or bans the trade of derived products. Despite being legally protected in many countries, the main causes of population decline remain anthropogenic and include fishing-related fatalities, habitat degradation, and hunting. With its long lifespan of 70 years or more and slow rate of reproduction, the dugong is especially vulnerable to extinction.
Evolution
Dugongs are part of the Sirenia order of placental mammals which comprises modern "sea cows" (manatees as well as dugongs) and their extinct relatives. Sirenia are the only extant herbivorous marine mammals and the only group of herbivorous mammals to have become completely aquatic. Sirenians are thought to have a 50-million-year-old fossil record (early Eocene-recent). They attained modest diversity during the Oligocene and Miocene but subsequently declined as a result of climatic cooling, oceanographic changes, and human interference.
Etymology and taxonomy
The word "dugong" derives from the Visayan (probably Cebuano) . The name was first adopted and popularized by the French naturalist Georges-Louis Leclerc, Comte de Buffon, as "dugon" in Histoire Naturelle (1765), after descriptions of the animal from the island of Leyte in the Philippines. The name ultimately derives from Proto-Malayo-Polynesian *duyuŋ. Despite common misconception, the term does not come from Malay duyung and it does not mean "lady of the sea" (Mermaid).
Other common local names include "sea cow", "sea pig" and "sea camel". It is known as the balguja by the Wunambal people of the Mitchell Plateau area in the Kimberley, Western Australia.
Dugong dugon is the only extant species of the family Dugongidae, and one of only four extant species of the Sirenia order, the others forming the manatee family. It was first classified by Müller in 1776 as Trichechus dugon, a member of the manatee genus previously defined by Linnaeus. It was later assigned as the type species of Dugong by Lacépède and further classified within its own family by Gray and subfamily by Simpson.
Dugongs and other sirenians are not closely related to other marine mammals, being more related to elephants. Dugongs and elephants share a monophyletic group with hyraxes and the aardvark, one of the earliest offshoots of eutherians. The fossil record shows sirenians appearing in the Eocene, where they most likely lived in the Tethys Ocean. The two extant families of sirenians are thought to have diverged in the mid-Eocene, after which the dugongs and their closest relative, the Steller's sea cow, split off from a common ancestor in the Miocene. The Steller's sea cow became extinct in the 18th century. No fossils exist of other members of the Dugongidae.
Molecular studies have been made on dugong populations using mitochondrial DNA. The results have suggested that the population of Southeast Asia is distinct from the others. Australia has two distinct maternal lineages, one of which also contains the dugongs from Africa and Arabia. Limited genetic mixing has taken place between those in Southeast Asia and those in Australia, mostly around Timor. One of the lineages stretches from Moreton Bay to Western Australia, while the other only stretches from Moreton Bay to the Northern Territory. There is not yet sufficient genetic data to make clear boundaries between distinct groups.
Anatomy and morphology
The dugong's body is large with a cylindrical shape that tapers at both ends. It has thick, smooth skin that is a pale cream colour at birth, but darkens dorsally and laterally to brownish-to-dark-grey with age. The colour of a dugong can change due to the growth of algae on the skin. The body is sparsely covered in short hair, a common feature among sirenians which may allow for tactile interpretation of their environment. These hairs are most developed around the mouth, which has a large horseshoe-shaped upper lip forming a highly mobile muzzle. This muscular upper lip aids the dugong in foraging.
The dugong's tail flukes and flippers are similar to those of dolphins. These flukes are raised up and down in long strokes to move the animal forward and can be twisted to turn. The forelimbs are paddle-like flippers which aid in turning and slowing. The dugong lacks nails on its flippers, which are only 15% of a dugong's body length. The tail has deep notches.
A dugong's brain weighs a maximum of , about 0.1% of the animal's body weight. With very small eyes, dugongs have limited vision, but acute hearing within narrow sound thresholds. Their ears, which lack pinnae, are located on the sides of their head. The nostrils are located on top of the head and can be closed using valves. Dugongs have two teats, one located behind each flipper. There are few differences between the sexes; the body structures are almost the same. A male's testes are not externally located, and the main difference between males and females is the location of the genital aperture to the umbilicus and the anus. The lungs in a dugong are very long, extending almost as far as the kidneys, which are also highly elongated to cope with the saltwater environment. If the dugong is wounded, its blood will clot rapidly.
The skull of a dugong is unique. The skull is enlarged with a sharply down-turned premaxilla, which is stronger in males. The spine has between 57 and 60 vertebrae. Unlike in manatees, the dugong's teeth do not continually grow back via horizontal tooth replacement. The dugong has two incisors (tusks) which emerge in males during puberty. The female's tusks continue to grow without emerging during puberty, sometimes erupting later in life after reaching the base of the premaxilla. The number of growth layer groups in a tusk indicates the age of a dugong, and the cheek teeth move forward with age.
The full dental formula of dugongs is , meaning they have two incisors, three premolars, and three molars on each side of their upper jaw, and three incisors, one canine, three premolars, and three molars on each side of their lower jaw. Like other sirenians, the dugong experiences pachyostosis, a condition in which the ribs and other long bones are unusually solid and contain little or no marrow. These heavy bones, which are among the densest in the animal kingdom, may act as a ballast to help keep sirenians suspended slightly below the water's surface.
An adult's length rarely exceeds . An individual this long is expected to weigh around . Weight in adults is typically more than and less than . The largest individual recorded was long and weighed , and was found off the Saurashtra coast of west India. Females tend to be larger than males.
Distribution and habitat
Dugongs are found in warm coastal waters from the western Pacific Ocean to the eastern coast of Africa, along an estimated of coastline between 26° and 27° to the north and south of the equator. Their historic range is believed to correspond to that of seagrasses from the Potamogetonaceae and Hydrocharitaceae families. The full size of the former range is unknown, although it is believed that the current populations represent the historical limits of the range, which is highly fractured. Their distributions during warmer periods of Holocene might have been broader than today. Today populations of dugongs are found in the waters of 37 countries and territories. Recorded numbers of dugongs are generally believed to be lower than actual numbers, due to a lack of accurate surveys. Despite this, the dugong population is thought to be shrinking, with a worldwide decline of 20 percent in the last 90 years. They have disappeared from the waters of Hong Kong, Mauritius, and Taiwan, as well as parts of Cambodia, Japan, the Philippines, and Vietnam. Further disappearances are likely.
Dugongs are generally found in warm waters around the coast with large numbers concentrated in wide and shallow protected bays. The dugong is the only strictly marine herbivorous mammal, as all species of manatee utilise fresh water to some degree. Nonetheless, they can tolerate the brackish waters found in coastal wetlands, and large numbers are also found in wide and shallow mangrove channels and around leeward sides of large inshore islands, where seagrass beds are common. They are usually located at a depth of around , although in areas where the continental shelf remains shallow dugongs have been known to travel more than from the shore, descending to as far as , where deepwater seagrasses such as Halophila spinulosa are found. Special habitats are used for different activities. It has been observed that shallow waters are used as sites for calving, minimizing the risk of predation. Deep waters may provide a thermal refuge from cooler waters closer to the shore during winter.
Australia
Australia is home to the largest population, stretching from Shark Bay in Western Australia to Moreton Bay in Queensland. The population of Shark Bay is thought to be stable with over 10,000 dugongs. Smaller populations exist up the coast, including one in Ashmore Reef. Large numbers of dugongs live to the north of the Northern Territory, with a population of over 20,000 in the Gulf of Carpentaria alone. A population of over 25,000 exists in the Torres Strait such as off Thursday Island, although there is significant migration between the strait and the waters of New Guinea.
The Great Barrier Reef provides important feeding areas for the species; this reef area houses a stable population of around 10,000, although the population concentration has shifted over time. Large bays facing north on the Queensland coast provide significant habitats for dugong, with the southernmost of these being Hervey Bay and Moreton Bay. Dugongs had been occasional visitors along the Gold Coast where a re-establishment of a local population through range expansions has started recently.
Persian Gulf
The Persian Gulf has the second-largest dugong population in the world, inhabiting most of the southern coast, and the current population is believed to range from 5,800 to 7,300. In the course of a study carried out in 1986 and 1999 on the Persian Gulf, the largest reported group sighting was made of more than 600 individuals to the west of Qatar. A 2017 study found a nearly 25% drop in population since 1950. Reasons for this drastic population loss include illegal poaching, oil spills, and net entanglement.
East Africa and South Asia
In the late 1960s, herds of up to 500 dugongs were observed off the coast of East Africa and nearby islands. Current populations in this area are extremely small, numbering 50 and below, and it is thought likely they will become extinct. The eastern side of the Red Sea is home to large populations numbering in the hundreds, and similar populations are thought to exist on the western side. In the 1980s, it was estimated there could be as many as 4,000 dugongs in the Red Sea. Dugong populations in Madagascar are poorly studied, but due to widespread exploitation, it is thought they may have severely declined, with few surviving individuals. The resident population around Mayotte is thought to number just 10 individuals. In Mozambique, most of the remaining local populations are very small and the largest (about 120 individuals) occurs at Bazaruto Island, but they have become rare in historical habitats such as in Maputo Bay and on Inhaca Island. The Bazaruto Island population is possibly the last long-term viable population in East Africa, with only some of its core territory lying within protected waters.
The East African population is genetically distinct from those of the Red Sea and those off Madagascar. In Tanzania, observations have recently increased around the Mafia Island Marine Park where a hunt was intended by fishermen but failed in 2009. In the Seychelles, dugongs had been regarded as extinct in the 18th century until a small number was discovered around the Aldabra Atoll. This population may belong to a different group than that distributed among the inner isles. Dugongs once thrived among the Chagos Archipelago and Sea Cow Island was named after the species, although the species no longer occurs in the region.
There are less than 250 individuals scattered throughout Indian waters. A highly isolated breeding population exists in the Marine National Park, Gulf of Kutch, the only remaining population in western India. It is from the population in the Persian Gulf, and from the nearest population in India. Former populations in this area, centered on the Maldives and the Lakshadweep, are presumed to be extinct. A population exists in the Gulf of Mannar Marine National Park and the Palk Strait between India and Sri Lanka, but it is seriously depleted. Recoveries of seagrass beds along former ranges of dugongs, such as the Chilika Lake have been confirmed in recent years, raising hopes for re-colorizations of the species. The population around the Andaman and Nicobar Islands is known only from a few records, and although the population was large during British rule, it is now believed to be small and scattered.
Southeast Asia and the West Pacific
A small population existed along the southern coast of China, particularly the Gulf of Tonkin (Beibu Gulf), where efforts were made to protect it, including the establishment of a seagrass sanctuary for dugong and other endangered marine fauna ranging in Guangxi. Despite these efforts, numbers continued to decrease, and in 2007 it was reported that no more dugong could be found on the west coast of the island of Hainan. Historically, dugongs were also present in the southern parts of the Yellow Sea. The last confirmed record of dugongs in Chinese waters was documented in 2008. In August 2022, an article published on the Royal Society Open Science concluded that dugongs were functionally extinct in China, which was based on a large-scale interview survey conducted across four southern Chinese maritime provinces (Hainan, Guangxi, Guangdong, and Fujian) in the summer of 2019.
In Vietnam, dugongs have been restricted mostly to the provinces of Kiên Giang and Bà Rịa–Vũng Tàu, including Phu Quoc Island and Con Dao Island, which hosted large populations in the past. Con Dao is now the only site in Vietnam where dugongs are regularly seen, protected within the Côn Đảo National Park. Nonetheless, dangerously low levels of attention to the conservation of marine organisms in Vietnam and Cambodia may result in increased intentional or unintentional catches, and illegal trade is a potential danger for local dugongs. On Phu Quoc, the first 'Dugong Festival' was held in 2014, aiming to raise awareness of these issues.
In Thailand, the present distribution of dugongs is restricted to six provinces along the Andaman Sea, and very few dugongs are present in the Gulf of Thailand. The Gulf of Thailand was historically home to a large number of animals, but none have been sighted in the west of the gulf in recent years, and the remaining population in the east is thought to be very small and possibly declining. Dugongs are believed to exist in the Straits of Johor in very small numbers. The waters around Borneo support a small population, with more scattered throughout the Malay Archipelago.
All the islands of the Philippines once provided habitats for sizeable herds of dugongs. They were common until the 1970s when their numbers declined sharply due to accidental drownings in fishing gear and habitat destruction of seagrass meadows. Today, only isolated populations survive, most notably in the waters of the Calamian Islands in Palawan, Isabela in Luzon, Guimaras, and Mindanao. The dugong became the first marine animal protected by Philippine law, with harsh penalties for harming them. Recently, the local marine trash problem in the archipelago remained unabated and became the biggest threat to the already dwindling population of Dugongs in the country. Litters of plastic waste (single-use sachets, plastic bottles, fast food to-go containers, etc.) and other non-biodegradable materials abound in the coastal areas. As these materials may be mistaken as food by dugongs, these may lead to death due to plastic ingestion. Overpopulation and lack of education of all coastal fisherfolk in the Philippines regarding marine trash are harming the coastal environment not only in Palawan but also across the islands of the Philippines. The first documented sighting in Sarangani Bay occurred in July 2024.
Populations also exist around the Solomon Islands and New Caledonia, stretching to an easternmost population in Vanuatu. A highly isolated population lives around the islands of Palau.
A single dugong lives at Cocos (Keeling) Islands although the animal is thought to be a vagrant.
Northern Pacific
Today, possibly the smallest and northernmost population of dugongs exists around the Ryukyu islands, and a population formerly existed off Taiwan. An endangered population of 50 or fewer dugongs, possibly as few as three individuals, survives around Okinawa. New sightings of a cow and calf have been reported in 2017, indicating a possible breeding had occurred in these waters. A single individual was recorded at Amami Ōshima, at the northernmost edge of the dugong's historic range, more than 40 years after the last previous recorded sighting. A vagrant strayed into a port near Ushibuka, Kumamoto, and died due to poor health. Historically, the Yaeyama Islands held a large concentration of dugongs, with more than 300 individuals. On the Aragusuku Islands, large quantities of skulls are preserved at a utaki that outsiders are strictly forbidden to enter. Dugong populations in these areas were reduced by historical hunts as payments to the Ryukyu Kingdom, before being wiped out because of large-scale illegal hunting and fishing using destructive methods such as dynamite fishing after the Second World War.
Populations around Taiwan appear to be almost extinct, although remnant individuals may visit areas with rich seagrass beds such as Dongsha Atoll. Some of the last reported sightings were made in Kenting National Park in the 1950s and 60s. There had been occasional records of vagrants at the Northern Mariana Islands before 1985. It is unknown how much mixing there was between these populations historically. Some theorize that populations existed independently, for example, that the Okinawan population was isolated members derived from the migration of a Philippine subspecies. Others postulate that the populations formed part of a super-population where migration between Ryukyu, Taiwan, and the Philippines was common.
Extinct Mediterranean population
It has been confirmed that dugongs once inhabited the water of the Mediterranean possibly until after the rise of civilizations along the inland sea. This population possibly shared ancestry with the Red Sea population, and the Mediterranean population had never been large due to geographical factors and climate changes. The Mediterranean is the region where the Dugongidae originated in the mid-late Eocene, along with Caribbean Sea.
Ecology and life history
Dugongs are long-lived, and the oldest recorded specimen reached age 73. They have few natural predators, although animals such as crocodiles, killer whales, and sharks pose a threat to the young, and a dugong has also been recorded to have died from trauma after being impaled by a stingray barb. A large number of infections and parasitic diseases affect dugongs. Detected pathogens include helminths, cryptosporidium, different types of bacterial infections, and other unidentified parasites. 30% of dugong deaths in Queensland since 1996 are thought to be because of disease.
Although they are social animals, they are usually solitary or found in pairs due to the inability of seagrass beds to support large populations. Gatherings of hundreds of dugongs sometimes happen, but they last only for a short time. Because they are shy and do not approach humans, little is known about dugong behavior. They can go six minutes without breathing (though about two and a half minutes is more typical), and have been known to rest on their tails to breathe with their heads above water. They can dive to a maximum depth of ; they spend most of their lives no deeper than . Communication between individuals is through chirps, whistles, barks, and other sounds that echo underwater. Different sounds have been observed with different amplitudes and frequencies, implying different purposes. Visual communication is limited due to poor eyesight and is mainly used for activities such as lekking for courtship purposes. Mothers and calves are in almost constant physical contact, and calves have been known to reach out and touch their mothers with their flippers for reassurance.
Dugongs are semi-nomadic, often traveling long distances in search of food, but staying within a certain range their entire lives. Large numbers often move together from one area to another. It is thought that these movements are caused by changes in seagrass availability. Their memory allows them to return to specific points after long travels. Dugong movements mostly occur within a localized area of seagrass beds, and animals in the same region show individualistic patterns of movement. Daily movement is affected by the tides. In areas where there is a large tidal range, dugongs travel with the tide to access shallower feeding areas. In Moreton Bay, dugongs often travel between foraging grounds inside the bay and warmer oceanic waters. At higher latitudes dugongs make seasonal travels to reach warmer water during the winter. Occasionally individual dugongs make long-distance travels over many days and can travel over deep ocean waters. One animal was seen as far south as Sydney. Although they are marine creatures, dugongs have been known to travel up creeks, and in one case a dugong was caught up a creek near Cooktown.
Feeding
Dugongs, along with other sirenians, are referred to as "sea cows" because their diet consists mainly of seagrass, particularly the genera Halophila and Halodule. When eating they ingest the whole plant, including the roots, although when this is impossible they will feed on just the leaves. A wide variety of seagrass has been found in dugong stomach contents, and evidence exists they will eat algae when seagrass is scarce. Although almost completely herbivorous, they will occasionally eat invertebrates such as jellyfish, sea squirts, and shellfish. Dugongs in Moreton Bay, Australia, are omnivorous, feeding on invertebrates such as polychaetes or marine algae when the supply of their choice grasses decreases. In other southern areas of both western and eastern Australia, there is evidence that dugongs actively seek out large invertebrates. This does not apply to dugongs in tropical areas, in which fecal evidence indicates that invertebrates are not eaten.
Most dugongs do not feed on lush areas, but where the seagrass is more sparse. Additional factors such as protein concentration and regenerative ability also affect the value of a seagrass bed. The chemical structure and composition of the seagrass are important, and the grass species most often eaten are low in fiber, high in nitrogen, and easily digestible. In the Great Barrier Reef, dugongs feed on low-fiber high-nitrogen seagrass such as Halophila and Halodule, to maximize nutrient intake instead of bulk eating. Seagrasses of a lower seral are preferred, where the area has not fully vegetated. Only certain seagrass meadows are suitable for dugong consumption, due to the dugong's highly specialized diet. There is evidence that dugongs actively alter seagrass species compositions at local levels. Dugongs may search out deeper seagrass. Feeding trails have been observed as deep as , and dugongs have been seen feeding as deep as . Dugongs are relatively slow-moving, swimming at around . When moving along the seabed to feed they walk on their pectoral fins.
Dugong feeding may favor the subsequent growth of low-fibre, high-nitrogen seagrasses such as Halophilia and Halodule. Species such as Zosteria capricorni are more dominant in established seagrass beds, but grow slowly, while Halophilia and Halodule grow quickly in the open space left by dugong feeding. This behavior is known as cultivation grazing and favors the rapidly growing, higher nutrient seagrasses that dugongs prefer. Dugongs may also prefer to feed on younger, less fibrous strands of seagrasses, and cycles of cultivation feeding at different seagrass meadows may provide them with a greater number of younger plants.
Due to their poor eyesight, dugongs often use smell to locate edible plants. They also have a strong tactile sense and feel their surroundings with their long sensitive bristles. They will dig up an entire plant and then shake it to remove the sand before eating it. They have been known to collect a pile of plants in one area before eating them. The flexible and muscular upper lip is used to dig out the plants. This leaves furrows in the sand in their path.
Reproduction and parental care
A dugong reaches sexual maturity between the ages of eight and eighteen, older than in most other mammals. The way that females know how a male has reached sexual maturity is by the eruption of tusks in the male since tusks erupt in males when testosterone levels reach a high enough level. The age when a female first gives birth is disputed, with some studies placing the age between ten and seventeen years, while others place it as early as six years. There is evidence that male dugongs lose fertility at older ages.
Despite the longevity of the dugong, which may live for 50 years or more, females give birth only a few times during their lives and invest considerable parental care in their young. The time between births is unclear, with estimates ranging from 2.4 to 7 years.
Mating behaviour varies between populations located in different areas. In some populations, males will establish a territory that females in estrus will visit. In these areas, a male will try to impress the females while defending the area from other males, a practice known as lekking. In other areas many males will attempt to mate with the same female, sometimes inflicting injuries to the female or each other. During this, the female will have copulated with multiple males, who will have fought to mount her from below. This greatly increases the chances of conception.
Females give birth after a 13- to 15-month gestation, usually to just one calf. Birth occurs in very shallow water, with occasions known where the mothers were almost on the shore. As soon as the young is born the mother pushes it to the surface to take a breath. Newborns are already long and weigh around . Once born, they stay close to their mothers, possibly to make swimming easier. The calf nurses for 14–18 months, although it begins to eat seagrasses soon after birth. A calf will only leave its mother once it has matured.
Importance to humans
Dugongs have historically provided easy targets for hunters, who killed them for their meat, oil, skin, and bones. As the anthropologist A. Asbjørn Jøn has noted, they are often considered the inspiration for mermaids, and people around the world developed cultures around dugong hunting. In some areas, it remains an animal of great significance, and a growing ecotourism industry around dugongs has had an economic benefit in some countries.
There is a 5,000-year-old wall painting of a dugong, apparently drawn by Neolithic peoples, in Tambun Cave, Ipoh, Malaysia. This was discovered by Lieutenant R.L. Rawlings in 1959 while on a routine patrol.
Dugongs feature in Southeast Asian, especially Austronesian, folklore. In languages like Ilocano, Mapun, Yakan, Tausug, and Kadazan Dusun of the Philippines and Sabah, the name for dugongs is a synonym for "mermaid". In Malay, they are sometimes referred to as perempoen laut ("woman of the sea") or putri duyong ("dugong princess"), leading to the misconception that the word "dugong" itself means "lady of the sea". A common belief found in the Philippines, Malaysia, Indonesia, and Thailand, is that dugongs were originally human or part-human (usually women), and that they cry when they are butchered or beached. Because of this, it is considered bad luck if a dugong is killed or accidentally dies in nets or fish corrals in the Philippines, some parts of Sabah (Malaysia), and northern Sulawesi and the Lesser Sunda Islands (Indonesia). Dugongs are predominantly not traditionally hunted for food in these regions and they remained plentiful until around the 1970s.
Conversely, dugong "tears" are considered aphrodisiacs in other parts of Indonesia, Singapore, Malaysia, Brunei, Thailand, Vietnam, and Cambodia. Dugong meat is considered a luxury food and is also believed to have aphrodisiac properties. They are actively hunted in these regions, in some places to near-extinction.
In Palau, dugongs were traditionally hunted with heavy spears from canoes. Although it is illegal and there is widespread disapproval of killing dugongs, poaching remains a major problem. Dugongs are also widely hunted in Papua New Guinea, the Solomon Islands, Vanuatu, and New Caledonia; where their meat and ornaments made from bones and tusks are highly prized in feasts and traditional rituals. However, hunting dugongs is considered taboo in some areas of Vanuatu. Dugong meat and oil have traditionally been some of the most valuable foods of Australian Aboriginals and Torres Strait Islanders. Some Aboriginals regard dugongs as part of their Aboriginality.
Local fishermen in Southern China traditionally revered dugongs and regarded them as "miraculous fish". They believed it was bad luck to catch them and they were plentiful in the region before the 1960s. Beginning in the 1950s, a wave of immigrants from other regions that do not hold these beliefs resulted in dugongs being hunted for food and traditional Chinese medicine. This led to a steep decline in dugong populations in the Gulf of Tonkin and the sea around Hainan Island. In Japan, dugongs have been traditionally hunted in the Ryukyu Islands since prehistoric times. Carved ribs of dugongs in the shape of butterflies (a psychopomp) are found throughout Okinawa. They were commonly hunted throughout Japan up until around the 1970s.
Dugongs have also played a role in legends in Kenya, and the animal is known there as the "Queen of the Sea". Body parts are used as food, medicine, and decorations. In the Gulf states, dugongs served not only as a source of food but their tusks were used as sword handles. Dugong oil is important as a preservative and conditioner for wooden boats to people around the Gulf of Kutch in India, who also believe the meat to be an aphrodisiac.
Conservation
Dugong numbers have decreased in recent times. For a population to remain stable, the mortality of adults cannot exceed 5% annually. The estimated percentage of females humans can kill without depleting the population is 1–2%. This number is reduced in areas where calving is minimal due to food shortages. Even in the best conditions, a population is unlikely to increase more than 5% a year, leaving dugongs vulnerable to over-exploitation. The fact that they live in shallow waters puts them under great pressure from human activity. Research on dugongs and the effects of human activity on them has been limited, mostly taking place in Australia. In many countries, dugong numbers have never been surveyed. As such, trends are uncertain, with more data needed for comprehensive management. The only data stretching back far enough to mention population trends comes from the urban coast of Queensland, Australia. The last major worldwide study, made in 2002, concluded that the dugong was declining and possibly extinct in a third of its range, with unknown status in another half.
The IUCN Red List lists the dugong as vulnerable, and the Convention on International Trade in Endangered Species of Wild Fauna and Flora regulates and in some areas has banned international trade. Most dugong habitats fall within proposed important marine mammal areas. Regional cooperation is important due to the widespread distribution of the animal, and in 1998 there was strong support for Southeast Asian cooperation to protect dugongs. Kenya has passed legislation banning the hunting of dugongs and restricting trawling, but the dugong is not yet listed under Kenya's Wildlife Act as an endangered species. Mozambique has had legislation to protect dugongs since 1955, but this has not been effectively enforced. France has a National Action Plan covering the species, implemented within the Mayotte Marine Natural Park.
Many marine parks have been established on the African coast of the Red Sea, and the Egyptian Gulf of Aqaba is fully protected. The United Arab Emirates has banned all hunting of dugongs within its waters, as has Bahrain. The UAE has additionally banned drift net fishing, and has declared an intention to restore coastal ecosystems dugongs rely on. India and Sri Lanka ban the hunting and selling of dugongs and their products. Japan has listed dugongs as endangered and has banned intentional killing and harassment. Hunting, catching, and harassment are banned by the People's Republic of China. The first marine mammal to be protected in the Philippines was the dugong, although monitoring this is difficult. Palau has legislated to protect dugongs, although this is not well enforced and poaching persists. Indonesia listed dugongs as a protected species in 1999, and in 2018 the Fisheries Ministry began implementing a conservation plan. Protection is not always enforced and souvenir products made from dugong parts can be openly found in markets in Bali. Traditional dugong hunters continued to hunt for many years, and some have struggled to find alternative incomes after ceasing. The dugong is a national animal of Papua New Guinea, which bans all except traditional hunting. Vanuatu and New Caledonia ban the hunting of dugongs. Dugongs are protected throughout Australia, although the rules vary by state; in some areas, indigenous hunting is allowed.
Dugongs are listed under the Nature Conservation Act in the Australian state of Queensland as vulnerable. Most currently live in established marine parks, where boats must travel at a restricted speed and mesh net fishing is restricted. The World Wide Fund for Nature has purchased gillnet licences in northern Queensland to reduce the impact of fishing. In Vietnam, an illegal network targeting dugongs had been detected and was shut down in 2012. Potential hunts along Tanzanian coasts by fishermen have raised concerns as well.
Human activity
Despite being legally protected in many countries, the main causes of population decline remain anthropogenic and include hunting, habitat degradation, and fishing-related fatalities. Entanglement in fishing nets has caused many deaths, although there are no precise statistics. Most issues with industrial fishing occur in deeper waters where dugong populations are low, with local fishing being the main risk in shallower waters. As dugongs cannot stay underwater for a very long period, they are highly prone to death due to entanglement. The use of shark nets has historically caused large numbers of deaths, and they have been eliminated in most areas and replaced with baited hooks. Hunting has historically been a problem too, although in most areas they are no longer hunted, except in certain indigenous communities. In areas such as northern Australia, hunting has the greatest impact on the dugong population.
Vessel strikes have proved a problem for manatees, but the relevance of this to dugongs is unknown. Increasing boat traffic has increased danger, especially in shallow waters. Ecotourism has increased in some countries, although the effects remain undocumented. It has been seen to cause issues in areas such as Hainan due to environmental degradation. Modern farming practices and increased land clearing have also had an impact, and much of the coastline of dugong habitats is undergoing industrialization, with increasing human populations. Dugongs accumulate heavy metal ions in their tissues throughout their lives, more so than other marine mammals. The effects are unknown.
While international cooperation to form a conservative unit has been undertaken, socio-political needs are an impediment to dugong conservation in many developing countries. The shallow waters are often used as a source of food and income, problems exacerbated by aid used to improve fishing. In many countries, legislation does not exist to protect dugongs, and if it does it is not enforced.
Oil spills are a danger to dugongs in some areas, as is land reclamation. In Okinawa, the small dugong population is threatened by United States military activity. Plans exist to build a military base close to the Henoko reef, and military activity also adds the threats of noise pollution, chemical pollution, soil erosion, and exposure to depleted uranium. The military base plans have been fought in US courts by some Okinawans, whose concerns include the impact on the local environment and dugong habitats. It was later revealed that the government of Japan was hiding evidence of the negative effects of ship lanes and human activities on dugongs observed during surveys carried out off Henoko reef. One of the three individuals has not been observed since June 2015, corresponding to the start of the excavation operations.
Environmental degradation
If dugongs do not get enough to eat they may calve later and produce fewer young. Food shortages can be caused by many factors, such as a loss of habitat, death and decline in the quality of seagrass, and a disturbance of feeding caused by human activity. Sewage, detergents, heavy metals, hypersaline water, herbicides, and other waste products all negatively affect seagrass meadows. Human activity such as mining, trawling, dredging, land reclamation, and boat propeller scarring also cause an increase in sedimentation which smothers seagrass and prevents light from reaching it. This is the most significant negative factor affecting seagrass.
Halophila ovalis—one of the dugong's preferred species of seagrass—declines rapidly due to lack of light, dying completely after 30 days. Extreme weather such as cyclones and floods can destroy hundreds of square kilometres of seagrass meadows, as well as wash dugongs ashore. The recovery of seagrass meadows and the spread of seagrass into new areas, or areas where it has been destroyed, can take over a decade. Most measures for protection involve restricting activities such as trawling in areas containing seagrass meadows, with little to no action on pollutants originating from land. In some areas, water salinity is increased due to wastewater, and it is unknown how much salinity seagrass can withstand.
Dugong habitat in the Oura Bay area of Henoko, Okinawa, Japan, is currently under threat from land reclamation conducted by the Japanese Government in order to build a US Marine base in the area. In August 2014, preliminary drilling surveys were conducted around the seagrass beds there. The construction is expected to seriously damage the dugong population's habitat, possibly leading to local extinction.
Capture and captivity
The Australian state of Queensland has sixteen dugong protection parks, and some preservation zones have been established where even Aboriginal Peoples are not allowed to hunt. Capturing animals for research has caused only one or two deaths; dugongs are expensive to keep in captivity due to the long time mothers and calves spend together, and the inability to grow the seagrass that dugongs eat in an aquarium. Only one orphaned calf has ever been successfully kept in captivity.
Worldwide, only three dugongs are held in captivity. A female from the Philippines lives at Toba Aquarium in Toba, Mie, Japan. A male also lived there until he died on 10 February 2011. The second resides in Sea World Indonesia, after having been rescued from a fisherman's net and treated. The last one, a male, is kept at Sydney Aquarium, where he has resided since he was a juvenile. Sydney Aquarium had a second dugong for many years, until she died in 2018.
Gracie, a captive dugong at Underwater World, Singapore, was reported to have died in 2014 at the age of 19, from complications arising from an acute digestive disorder.
| Biology and health sciences | Sirenia | Animals |
60257 | https://en.wikipedia.org/wiki/Marine%20mammal | Marine mammal | Marine mammals are mammals that rely on marine (saltwater) ecosystems for their existence. They include animals such as cetaceans (whales, dolphins and porpoises), pinnipeds (seals, sea lions and walruses), sirenians (manatees and dugongs), sea otters and polar bears. They are an informal group, unified only by their reliance on marine environments for feeding and survival.
Marine mammal adaptation to an aquatic lifestyle varies considerably between species. Both cetaceans and sirenians are fully aquatic and therefore are obligate water dwellers. Pinnipeds are semiaquatic; they spend the majority of their time in the water but need to return to land for important activities such as mating, breeding and molting. Sea otters tend to live in kelp forests and estuaries. In contrast, the polar bear is mostly terrestrial and only go into the water on occasions of necessity, and are thus much less adapted to aquatic living. The diets of marine mammals vary considerably as well; some eat zooplankton, others eat fish, squid, shellfish, or seagrass, and a few eat other mammals. While the number of marine mammals is small compared to those found on land, their roles in various ecosystems are large, especially concerning the maintenance of marine ecosystems, through processes including the regulation of prey populations. This role in maintaining ecosystems makes them of particular concern as 23% of marine mammal species are currently threatened.
Marine mammals were first hunted by aboriginal peoples for food and other resources. Many were also the target for commercial industry, leading to a sharp decline in all populations of exploited species, such as whales and seals. Commercial hunting led to the extinction of the Steller's sea cow, sea mink, Japanese sea lion and Caribbean monk seal. After commercial hunting ended, some species, such as the gray whale and northern elephant seal, have rebounded in numbers; conversely, other species, such as the North Atlantic right whale, are critically endangered. Other than being hunted, marine mammals can be killed as bycatch from fisheries, where for example they can become entangled in nets and drown or starve. Increased ocean traffic causes collisions between fast ocean vessels and large marine mammals. Habitat degradation also threatens marine mammals and their ability to find and catch food. Noise pollution, for example, may adversely affect echolocating mammals, and the ongoing effects of global warming degrade Arctic environments.
Taxonomy
Classification of extant species
Order Cetartiodactyla
Suborder Whippomorpha
Family Balaenidae (right and bowhead whales), two genera and four species
Family Cetotheriidae (pygmy right whale), one species
Family Balaenopteridae (rorquals), two genera and eight species
Family Eschrichtiidae (gray whale), one species
Family Physeteridae (sperm whale), one species
Family Kogiidae (pygmy and dwarf sperm whales), one genus and two species
Family Monodontidae (narwhal and beluga), two genera and two species
Family Ziphiidae (beaked whales), six genera and 21 species
Family Delphinidae (oceanic dolphins), 17 genera and 38 species
Family Phocoenidae (porpoises), two genera and seven species
Order Sirenia (sea cows)
Family Trichechidae (manatees), three species
Family Dugongidae (dugongs), one species
Order Carnivora (carnivorans)
Suborder Caniformia
Family Mustelidae, two species
Family Ursidae (bears), one species
Infraorder Pinnipedia (sealions, walruses, seals)
Family Otariidae (eared seals), seven genera and 15 species
Family Odobenidae (walrus), one species
Family Phocidae (earless seals), 14 genera and 18 species
The term "marine mammal" encompasses all mammals whose survival depends entirely or almost entirely on the oceans, which have also evolved several specialized aquatic traits. In addition to the above, several other mammals have a great dependency on the sea without having become so anatomically specialized, otherwise known as "quasi-marine mammals". This term can include: the greater bulldog bat (Noctilio leporinus), the fish-eating bat (Myotis vivesi), the arctic fox (Vulpes lagopus) which often scavenges polar bear kills, coastal gray wolf (Canis lupus) populations which predominantly eat salmon and marine carcasses, the North Ronaldsay sheep (Ovis aries) which normally eats seaweed outside the lambing season, the Eurasian otter (Lutra lutra) which is usually found in freshwater but can be found along coastal Scotland, and others.
Evolution
Marine mammals form a diverse group of 129 species that rely on the ocean for their existence. They are an informal group unified only by their reliance on marine environments for feeding. Despite the diversity in anatomy seen between groups, improved foraging efficiency has been the main driver in their evolution. The level of dependence on the marine environment varies considerably with species. For example, dolphins and whales are completely dependent on the marine environment for all stages of their life; seals feed in the ocean but breed on land; and polar bears must feed on land.
The cetaceans became aquatic around 50 million years ago (mya). Based on molecular and morphological research, the cetaceans genetically and morphologically fall firmly within the Artiodactyla (even-toed ungulates). The term "Cetartiodactyla" reflects the idea that whales evolved within the ungulates. The term was coined by merging the name for the two orders, Cetacea and Artiodactyla, into a single word. Under this definition, the closest living land relative of the whales and dolphins is thought to be the hippopotamuses.
Sirenians, the sea cows, became aquatic around 40 million years ago. The first appearance of sirenians in the fossil record was during the early Eocene, and by the late Eocene, sirenians had significantly diversified. Inhabitants of rivers, estuaries, and nearshore marine waters, they were able to spread rapidly. The most primitive sirenian, †Prorastomus, was found in Jamaica, unlike other marine mammals which originated from the Old World (such as cetaceans). The first known quadrupedal sirenian was †Pezosiren from the early middle Eocene. The earliest known sea cows, of the families †Prorastomidae and †Protosirenidae, were both confined to the Eocene, and were pig-sized, four-legged, amphibious creatures. The first members of Dugongidae appeared by the middle Eocene. At this point, sea cows were fully aquatic.
Pinnipeds split from other caniforms 50 mya during the Eocene. Their evolutionary link to terrestrial mammals was unknown until the 2007 discovery of †Puijila darwini in early Miocene deposits in Nunavut, Canada. Like a modern otter, †Puijila had a long tail, short limbs and webbed feet instead of flippers. The lineages of Otariidae (eared seals) and Odobenidae (walrus) split almost 28 mya. Phocids (earless seals) are known to have existed for at least 15 mya, and molecular evidence supports a divergence of the Monachinae (monk seals) and Phocinae lineages 22 mya.
Fossil evidence indicates the sea otter (Enhydra) lineage became isolated in the North Pacific approximately two mya, giving rise to the now-extinct †Enhydra macrodonta and the modern sea otter, Enhydra lutris. The sea otter evolved initially in northern Hokkaidō and Russia, and then spread east to the Aleutian Islands, mainland Alaska, and down the North American coast. In comparison to cetaceans, sirenians, and pinnipeds, which entered the water approximately 50, 40, and 20 mya, respectively, the sea otter is a relative newcomer to marine life. In some respects though, the sea otter is more fully adapted to water than pinnipeds, which must haul out on land or ice to give birth.
Polar bears are thought to have diverged from a population of brown bears, Ursus arctos, that became isolated during a period of glaciation in the Pleistocene or from the eastern part of Siberia, (from Kamchatka and the Kolym Peninsula). The oldest known polar bear fossil is a 130,000-to-110,000-year-old jaw bone, found on Prince Charles Foreland in 2004. The mitochondrial DNA (mtDNA) of the polar bear diverged from the brown bear roughly 150,000 years ago. Further, some clades of brown bear, as assessed by their mtDNA, are more closely related to polar bears than to other brown bears, meaning that the polar bear might not be considered a species under some species concepts.
In general, terrestrial amniote invasions of the sea have become more frequent in the Cenozoic than they were in the Mesozoic. Factors contributing to this trend include the increasing productivity of near-shore marine environments, and the role of endothermy in facilitating this transition.
Distribution and habitat
Marine mammals are widely distributed throughout the globe, but their distribution is patchy and coincides with the productivity of the oceans. Species richness peaks at around 40° latitude, both north and south. This corresponds to the highest levels of primary production around North and South America, Africa, Asia and Australia. Total species range is highly variable for marine mammal species. On average most marine mammals have ranges which are equivalent or smaller than one-fifth of the Indian Ocean. The variation observed in range size is a result of the different ecological requirements of each species and their ability to cope with a broad range of environmental conditions. The high degree of overlap between marine mammal species richness and areas of human impact on the environment is of concern.
Most marine mammals, such as seals and sea otters, inhabit the coast. Seals, however, also use a number of terrestrial habitats, both continental and island. In temperate and tropical areas, they haul-out on to sandy and pebble beaches, rocky shores, shoals, mud flats, tide pools and in sea caves. Some species also rest on man-made structures, like piers, jetties, buoys and oil platforms. Seals may move further inland and rest in sand dunes or vegetation, and may even climb cliffs. Most cetaceans live in the open ocean, and species like the sperm whale may dive to depths of in search of food. Sirenians live in shallow coastal waters, usually living below sea level. However, they have been known to dive to to forage deep-water seagrasses. Sea otters live in protected areas, such as rocky shores, kelp forests, and barrier reefs, although they may reside among drift ice or in sandy, muddy, or silty areas.
Many marine mammals seasonally migrate. Annual ice contains areas of water that appear and disappear throughout the year as the weather changes, and seals migrate in response to these changes. In turn, polar bears must follow their prey. In Hudson Bay, James Bay, and some other areas, the ice melts completely each summer (an event often referred to as "ice-floe breakup"), forcing polar bears to go onto land and wait through the months until the next freeze-up. In the Chukchi and Beaufort seas, polar bears retreat each summer to the ice further north that remains frozen year-round. Seals may also migrate to other environmental changes, such as El Niño, and traveling seals may use various features of their environment to reach their destination including geomagnetic fields, water and wind currents, the position of the sun and moon and the taste and temperature of the water. Baleen whales famously migrate very long distances into tropical waters to give birth and raise young, possibly to prevent predation by killer whales. The gray whale has the longest recorded migration of any mammal, with one traveling from the Sea of Okhotsk to the Baja Peninsula. During the winter, manatees living at the northern end of their range migrate to warmer waters.
Adaptations
Marine mammals have a number of physiological and anatomical features to overcome the unique challenges associated with aquatic living. Some of these features are very species-specific. Marine mammals have developed a number of features for efficient locomotion such as torpedo-shaped bodies to reduce drag; modified limbs for propulsion and steering; tail flukes and dorsal fins for propulsion and balance. Marine mammals are adept at thermoregulation using dense fur or blubber, circulatory adjustments (counter-current heat exchange); and reduced appendages, and large size to prevent heat loss.
Marine mammals are able to dive for long periods. Both pinnipeds and cetaceans have large and complex blood vessel systems pushing large volumes of blood rich in myoglobin and hemoglobin, which serve to store greater quantities of oxygen. Other important reservoirs include muscles and the spleen which all have the capacity to hold a high concentration of oxygen. They are also capable of bradycardia (reduced heart rate), and vasoconstriction (shunting most of the oxygen to vital organs such as the brain and heart) to allow extended diving times and cope with oxygen deprivation. If oxygen is depleted (hypoxia), marine mammals can access substantial reservoirs of glycogen that support anaerobic glycolysis.
Sound travels differently through water, and therefore marine mammals have developed adaptations to ensure effective communication, prey capture, and predator detection. The most notable adaptation is the development of echolocation in whales and dolphins. Toothed whales emit a focused beam of high-frequency clicks in the direction that their head is pointing. Sounds are generated by passing air from the bony nares through the phonic lips. These sounds are reflected by the dense concave bone of the cranium and an air sac at its base. The focused beam is modulated by a large fatty organ known as the 'melon'. This acts like an acoustic lens because it is composed of lipids of differing densities.
Marine mammals have evolved a wide variety of features for feeding, which are mainly seen in their dentition. For example, the cheek teeth of pinnipeds and odontocetes are specifically adapted to capture fish and squid. In contrast, baleen whales have evolved baleen plates to filter feed plankton and small fish from the water.
Polar bears, otters, and fur seals have long, oily, and waterproof fur in order to trap air to provide insulation. In contrast, other marine mammals—such as whales, dolphins, porpoises, manatees, dugongs, and walruses—have lost long fur in favor of a thick, dense epidermis and a thickened fat layer (blubber) to prevent drag. Wading and bottom-feeding animals (such as manatees) need to be heavier than water in order to keep contact with the floor or to stay submerged. Surface-living animals (such as sea otters) need the opposite, and free-swimming animals living in open waters (such as dolphins) need to be neutrally buoyant in order to be able to swim up and down the water column. Typically, thick and dense bone is found in bottom feeders and low bone density is associated with mammals living in deep water. Some marine mammals, such as polar bears and otters, have retained four weight-bearing limbs and can walk on land like fully terrestrial animals.
Ecology
Dietary
All cetaceans are carnivorous and predatory. Toothed whales mostly feed on fish and cephalopods, followed by crustaceans and bivalves. Some may forage with other kinds of animals, such as other species of whales or certain species of pinnipeds. One common feeding method is herding, where a pod squeezes a school of fish into a small volume, known as a bait ball. Individual members then take turns plowing through the ball, feeding on the stunned fish. Coralling is a method where dolphins chase fish into shallow water to catch them more easily. Killer whales and bottlenose dolphins have also been known to drive their prey onto a beach to feed on it. Killer whales have been known to paralyze great white sharks and other sharks and rays by flipping them upside down. Other whales with a blunt snout and reduced dentition rely on suction feeding. Though carnivorous, they house gut flora similar to that of terrestrial herbivores, probably a remnant of their herbivorous ancestry.
Baleen whales use their baleen plates to sieve plankton, among others, out of the water; there are two types of methods: lunge-feeding and gulp-feeding. Lunge-feeders expand the volume of their jaw to a volume bigger than the original volume of the whale itself by inflating their mouth. This causes grooves on their throat to expand, increasing the amount of water the mouth can store. They ram a baitball at high speeds in order to feed, but this is only energy-effective when used against a large baitball. Gulp-feeders swim with an open mouth, filling it with water and prey. Prey must occur in sufficient numbers to trigger the whale's interest, be within a certain size range so that the baleen plates can filter it, and be slow enough so that it cannot escape.
Otters are the only marine animals that are capable of lifting and turning over rocks, which they often do with their front paws when searching for prey. The sea otter may pluck snails and other organisms from kelp and dig deep into underwater mud for clams. It is the only marine mammal that catches fish with its forepaws rather than with its teeth. Under each foreleg, sea otters have a loose pouch of skin that extends across the chest which they use to store collected food to bring to the surface. This pouch also holds a rock that is used to break open shellfish and clams, an example of tool use. The sea otters eat while floating on their backs, using their forepaws to tear food apart and bring to their mouths. Marine otters mainly feed on crustaceans and fish.
Pinnipeds mostly feed on fish and cephalopods, followed by crustaceans and bivalves, and then zooplankton and warm-blooded prey (like sea birds). Most species are generalist feeders, but a few are specialists. They typically hunt non-schooling fish, slow-moving or immobile invertebrates or endothermic prey when in groups. Solitary foraging species usually exploit coastal waters, bays and rivers. When large schools of fish or squid are available, pinnipeds hunt cooperatively in large groups, locating and herding their prey. Some species, such as California and South American sea lions, may forage with cetaceans and sea birds.
The polar bear is the most carnivorous species of bear, and its diet primarily consists of ringed (Pusa hispida) and bearded (Erignathus barbatus) seals. Polar bears hunt primarily at the interface between ice, water, and air; they only rarely catch seals on land or in open water. The polar bear's most common hunting method is still-hunting: The bear locates a seal breathing hole using its sense of smell, and crouches nearby for a seal to appear. When the seal exhales, the bear smells its breath, reaches into the hole with a forepaw, and drags it out onto the ice. The polar bear also hunts by stalking seals resting on the ice. Upon spotting a seal, it walks to within , and then crouches. If the seal does not notice, the bear creeps to within of the seal and then suddenly rushes to attack. A third hunting method is to raid the birth lairs that female seals create in the snow. They may also feed on fish.
Sirenians are referred to as "sea cows" because their diet consists mainly of seagrass. When eating, they ingest the whole plant, including the roots, although when this is impossible they feed on just the leaves. A wide variety of seagrass has been found in dugong stomach contents, and evidence exists they will eat algae when seagrass is scarce. West Indian manatees eat up to 60 different species of plants, as well as fish and small invertebrates to a lesser extent.
Keystone species
Sea otters are a classic example of a keystone species; their presence affects the ecosystem more profoundly than their size and numbers would suggest. They keep the population of certain benthic (sea floor) herbivores, particularly sea urchins, in check. Sea urchins graze on the lower stems of kelp, causing the kelp to drift away and die. Loss of the habitat and nutrients provided by kelp forests leads to profound cascade effects on the marine ecosystem. North Pacific areas that do not have sea otters often turn into urchin barrens, with abundant sea urchins and no kelp forest. Reintroduction of sea otters to British Columbia has led to a dramatic improvement in the health of coastal ecosystems, and similar changes have been observed as sea otter populations recovered in the Aleutian and Commander Islands and the Big Sur coast of California. However, some kelp forest ecosystems in California have also thrived without sea otters, with sea urchin populations apparently controlled by other factors. The role of sea otters in maintaining kelp forests has been observed to be more important in areas of open coast than in more protected bays and estuaries.
An apex predator affects prey population dynamics and defense tactics (such as camouflage). The polar bear is the apex predator within its range. Several animal species, particularly Arctic foxes (Vulpes lagopus) and glaucous gulls (Larus hyperboreus), routinely scavenge polar bear kills. The relationship between ringed seals and polar bears is so close that the abundance of ringed seals in some areas appears to regulate the density of polar bears, while polar bear predation in turn regulates density and reproductive success of ringed seals. The evolutionary pressure of polar bear predation on seals probably accounts for some significant differences between Arctic and Antarctic seals. Compared to the Antarctic, where there is no major surface predator, Arctic seals use more breathing holes per individual, appear more restless when hauled out on the ice, and rarely defecate on the ice. The fur of Arctic pups is white, presumably to provide camouflage from predators, whereas Antarctic pups all have dark fur.
Killer whales are apex predators throughout their global distribution, and can have a profound effect on the behavior and population of prey species. Their diet is very broad and they can feed on many vertebrates in the ocean including salmon, rays, sharks (even white sharks), large baleen whales, and nearly 20 species of pinniped. The predation of whale calves may be responsible for annual whale migrations to calving grounds in more tropical waters, where the population of killer whales is much lower than in polar waters. Prior to whaling, it is thought that great whales were a major food source; however, after their sharp decline, killer whales have since expanded their diet, leading to the decline of smaller marine mammals. A decline in Aleutian Islands sea otter populations in the 1990s was controversially attributed by some scientists to killer whale predation, although with no direct evidence. The decline of sea otters followed a decline in harbor seal and Steller sea lion populations, the killer whale's preferred prey, which in turn may be substitutes for their original prey, now reduced by industrial whaling.
Whale pump
A 2010 study considered whales to be a positive influence to the productivity of ocean fisheries, in what has been termed a "whale pump". Whales carry nutrients such as nitrogen from the depths back to the surface. This functions as an upward biological pump, reversing an earlier presumption that whales accelerate the loss of nutrients to the bottom. This nitrogen input in the Gulf of Maine is more than the input of all rivers combined emptying into the gulf, some each year.
Whales defecate at the ocean's surface; their excrement is important for fisheries because it is rich in iron and nitrogen. The whale feces are liquid and instead of sinking, they stay at the surface where phytoplankton feed off it.
Upon death, whale carcasses fall to the deep ocean and provide a substantial habitat for marine life. Evidence of whale falls in present-day and fossil records shows that deep-sea whale falls support a rich assemblage of creatures, with a global diversity of 407 species, comparable to other neritic biodiversity hotspots, such as cold seeps and hydrothermal vents. Deterioration of whale carcasses happens through a series of three stages. Initially, moving organisms, such as sharks and hagfish, scavenge soft tissue at a rapid rate over a period of months to as long as two years. This is followed by the colonization of bones and surrounding sediments (which contain organic matter) by enrichment opportunists, such as crustaceans and polychaetes, throughout a period of years. Finally, sulfophilic bacteria reduce the bones releasing hydrogen sulphide enabling the growth of chemoautotrophic organisms, which in turn, support other organisms such as mussels, clams, limpets, and sea snails. This stage may last for decades and supports a rich assemblage of species, averaging 185 species per site.
Interactions with humans
Threats
Due to the difficulty to survey populations, 38% of marine mammals are data deficient, especially around the Antarctic Polar Front. In particular, declines in the populations of completely marine mammals tend to go unnoticed 70% of the time.
Exploitation
Marine mammals were hunted by coastal aboriginal humans historically for food and other resources. These subsistence hunts still occur in Canada, Greenland, Indonesia, Russia, the United States, and several nations in the Caribbean. The effects of these are only localized, as hunting efforts were on a relatively small scale. Commercial hunting took this to a much greater scale and marine mammals were heavily exploited. This led to the extinction of the Steller's sea cow (Hydrodamalis gigas), sea mink (Neogale macrodon), Japanese sea lion (Zalophus japonicus), and the Caribbean monk seal (Neomonachus tropicalis). Today, populations of species that were historically hunted, such as blue whales (Balaenoptera musculus) and the North Pacific right whale (Eubalaena japonica), are much lower than their pre-whaling levels. Because whales generally have slow growth rates, are slow to reach sexual maturity, and have a low reproductive output, population recovery has been very slow.
A number of whales are still subject to direct hunting, despite the 1986 moratorium on commercial whaling set under the terms of the International Whaling Commission (IWC). There are only two nations remaining which sanction commercial whaling: Norway, where several hundred common minke whales are harvested each year; and Iceland, where quotas of 150 fin whales and 100 minke whales per year are set. Japan also harvests several hundred Antarctic and North Pacific minke whales each year, ostensibly for scientific research in accordance with the moratorium. However, the illegal trade of whale and dolphin meat is a significant market in Japan and some countries.
The most profitable furs in the fur trade were those of sea otters, especially the northern sea otter which inhabited the coastal waters between the Columbia River to the south and Cook Inlet to the north. The fur of the Californian southern sea otter was less highly prized and thus less profitable. After the northern sea otter was hunted to local extinction, maritime fur traders shifted to California until the southern sea otter was likewise nearly extinct. The British and American maritime fur traders took their furs to the Chinese port of Guangzhou (Canton), where they worked within the established Canton System. Furs from Russian America were mostly sold to China via the Mongolian trading town of Kyakhta, which had been opened to Russian trade by the 1727 Treaty of Kyakhta.
Commercial sealing was historically just as important as the whaling industry. Exploited species included harp seals, hooded seals, Caspian seals, elephant seals, walruses and all species of fur seal. The scale of seal harvesting decreased substantially after the 1960s, after the Canadian government reduced the length of the hunting season and implemented measures to protect adult females. Several species that were commercially exploited have rebounded in numbers; for example, Antarctic fur seals may be as numerous as they were prior to harvesting. The northern elephant seal was hunted to near extinction in the late 19th century, with only a small population remaining on Guadalupe Island. It has since recolonized much of its historic range, but has a population bottleneck. Conversely, the Mediterranean monk seal was extirpated from much of its former range, which stretched from the Mediterranean to the Black Sea and northwest Africa, and remains only in the northeastern Mediterranean and some parts of northwest Africa.
Polar bears can be hunted for sport in Canada with a special permit and accompaniment by a local guide. This can be an important source of income for small communities, as guided hunts bring in more income than selling the polar bear hide on markets. The United States, Russia, Norway, Greenland, and Canada allow subsistence hunting, and Canada distributes hunting permits to indigenous communities. The selling of these permits is a main source of income for many of these communities. Their hides can be used for subsistence purposes, kept as hunting trophies, or can be bought in markets.
Ocean traffic and fisheries
By-catch is the incidental capture of non-target species in fisheries. Fixed and drift gill nets cause the highest mortality levels for both cetaceans and pinnipeds, however, entanglements in long lines, mid-water trawls, and both trap and pot lines are also common. Tuna seines are particularly problematic for entanglement by dolphins. By-catch affects all cetaceans, both small and big, in all habitat types. However, smaller cetaceans and pinnipeds are most vulnerable as their size means that escape once they are entangled is highly unlikely and they frequently drown. While larger cetaceans are capable of dragging nets with them, the nets sometimes remain tightly attached to the individual and can impede the animal from feeding sometimes leading to starvation. Abandoned or lost nets and lines cause mortality through ingestion or entanglement. Marine mammals also get entangled in aquaculture nets, however, these are rare events and not prevalent enough to impact populations.
Vessel strikes cause death for a number of marine mammals, especially whales. In particular, fast commercial vessels such as container ships can cause major injuries or death when they collide with marine mammals. Collisions occur both with large commercial vessels and recreational boats and cause injury to whales or smaller cetaceans. The critically endangered North Atlantic right whale is particularly affected by vessel strikes. Tourism boats designed for whale and dolphin watching can also negatively impact on marine mammals by interfering with their natural behavior.
The fishery industry not only threatens marine mammals through by-catch, but also through competition for food. Large-scale fisheries have led to the depletion of fish stocks that are important prey species for marine mammals. Pinnipeds have been especially affected by the direct loss of food supplies and in some cases the harvesting of fish has led to food shortages or dietary deficiencies, starvation of young, and reduced recruitment into the population. As the fish stocks have been depleted, the competition between marine mammals and fisheries has sometimes led to conflict. Large-scale culling of populations of marine mammals by commercial fishers has been initiated in a number of areas in order to protect fish stocks for human consumption.
Shellfish aquaculture takes up space so in effect creates competition for space. However, there is little direct competition for aquaculture shellfish harvest. On the other hand, marine mammals regularly take finfish from farms, which creates significant problems for marine farmers. While there are usually legal mechanisms designed to deter marine mammals, such as anti-predator nets or harassment devices, individuals are often illegally shot.
Habitat loss and degradation
Habitat degradation is caused by a number of human activities. Marine mammals that live in coastal environments are the most likely to be affected by habitat degradation and loss. Developments such as sewage marine outfalls, moorings, dredging, blasting, dumping, port construction, hydroelectric projects, and aquaculture both degrade the environment and take up valuable habitat. For example, extensive shellfish aquaculture takes up valuable space used by coastal marine mammals for important activities such as breeding, foraging and resting.
Contaminants that are discharged into the marine environment accumulate in the bodies of marine mammals when they are stored unintentionally in their blubber along with energy. Contaminants that are found in the tissues of marine mammals include heavy metals, such as mercury and lead, but also organochlorides and polycyclic aromatic hydrocarbons. For example, these can cause disruptive effects on endocrine systems; impair the reproductive system, and lower the immune system of individuals, leading to a higher number of deaths. Other pollutants such as oil, plastic debris and sewage threaten the livelihood of marine mammals.
Noise pollution from anthropogenic activities is another major concern for marine mammals. This is a problem because underwater noise pollution interferes with the abilities of some marine mammals to communicate, and locate both predators and prey. Underwater explosions are used for a variety of purposes including military activities, construction and oceanographic or geophysical research. They can cause injuries such as hemorrhaging of the lungs, and contusion and ulceration of the gastrointestinal tract. Underwater noise is generated from shipping, the oil and gas industry, research, and military use of sonar and oceanographic acoustic experimentation. Acoustic harassment devices and acoustic deterrent devices used by aquaculture facilities to scare away marine mammals emit loud and noxious underwater sounds.
Two changes to the global atmosphere due to anthropogenic activity threaten marine mammals. The first is increases in ultraviolet radiation due to ozone depletion, and this mainly affects the Antarctic and other areas of the Southern Hemisphere. An increase in ultraviolet radiation has the capacity to decrease phytoplankton abundance, which forms the basis of the food chain in the ocean. The second effect of global climate change is global warming due to increased carbon dioxide levels in the atmosphere. Raised sea levels, rising sea temperatures and changed currents are expected to affect marine mammals by altering the distribution of important prey species, and changing the suitability of breeding sites and migratory routes. The Arctic food chain would be disrupted by the near extinction or migration of polar bears. Arctic sea ice is the polar bear's habitat. It has been declining at a rate of 13% per decade because the temperature is rising at twice the rate of the rest of the world. By the year 2050, up to two-thirds of the world's polar bears may vanish if the sea ice continues to melt at its current rate.
A study by evolutionary biologists at the University of Pittsburgh showed that the ancestors of many marine mammals stopped producing a certain enzyme that today protects against some neurotoxic chemicals called organophosphates, including those found in the widely used pesticides chlorpyrifos and diazinon. Marine mammals may be increasingly exposed to these compounds due to agricultural runoff reaching the world's oceans.
Protection
The Marine Mammal Protection Act of 1972 (MMPA) was passed on October 21, 1972, under president Richard Nixon to prevent the further depletion and possible extinction of marine mammal stocks. It prohibits the taking ("the act of hunting, killing, capture, and/or harassment of any marine mammal; or, the attempt at such") of any marine mammal without a permit issued by the Secretary. Authority to manage the MMPA was divided between the Secretary of the Interior through the U.S. Fish and Wildlife Service (Service), and the Secretary of Commerce, which is delegated to the National Oceanic and Atmospheric Administration (NOAA). The Marine Mammal Commission (MMC) was established to review existing policies and make recommendations to the Service and NOAA to better implement the MMPA. The Service is responsible for ensuring the protection of sea otters and marine otters, walruses, polar bears, the three species of manatees, and dugongs; and NOAA was given responsibility to conserve and manage pinnipeds (excluding walruses) and cetaceans.
The Act was updated on 1 January 2016 with a clause banning "the import of fish from fisheries that cannot prove they meet US standards for protecting marine mammals". The requirement to show that protection standards are met is hoped to compel countries exporting fish to the US to more strictly control their fisheries that no protected marine mammals are adversely affected by fishing.
The 1979 Convention on the Conservation of Migratory Species of Wild Animals (CMS) is the only global organization that conserves a broad range of animals, which includes marine mammals. Of the agreements made, three of them deal with the conservation of marine mammals: ACCOBAMS, ASCOBANS and the Wadden Sea Agreement. In 1982, the United Nations Convention on the Law of the Sea (UNCLOS) adopted a pollution prevention approach to conservation, which many other conventions at the time also adopted.
The Agreement on the Conservation of Cetaceans in the Black Sea, Mediterranean Sea and contiguous Atlantic area (ACCOBAMS), founded in 1996, specifically protects cetaceans in the Mediterranean area, and "maintains a favorable status", a direct action against whaling. There are 23 member states.
The Agreement on the Conservation of Small Cetaceans of the Baltic and North Seas (ASCOBANS) was adopted alongside ACCOBAMS to establish a special protection area for Europe's increasingly threatened cetaceans. Other anti-whaling efforts include a ten-year moratorium in 1986 by the IWC on all whaling, and an environmental agreement (a type of international law) the International Convention for the Regulation of Whaling which controlled commercial, scientific and subsistence whaling.
The Agreement on the Conservation of Seals in the Wadden Sea, enforced in 1991, prohibits the killing or harassment of seals in the Wadden Sea, specifically targeting the harbor seal population.
The 1973 Agreement on the Conservation of Polar Bears between Canada, Denmark (Greenland), Norway (Svalbard), the United States and the Soviet Union outlawed the unregulated hunting of polar bears from aircraft and icebreakers, as well as protecting migration, feeding and hibernation sites.
Various non-governmental organizations participate in marine conservation activism, wherein they draw attention to and aid in various problems in marine conservation, such as pollution, whaling, bycatch, and so forth. Notable organizations include the Greenpeace who focus on overfishing and whaling among other things, and Sea Shepherd Conservation Society who are known for taking direct-action tactics to expose illegal activity.
As food
For thousands of years, indigenous peoples of the Arctic have depended on whale meat and seal meat. The meat is harvested from legal, non-commercial hunts that occur twice a year in the spring and autumn. The meat is stored and eaten throughout the winter. The skin and blubber (muktuk) taken from the bowhead, beluga, or narwhal is also valued, and is eaten raw or cooked. Whaling has also been practiced in the Faroe Islands in the North Atlantic since about the time of the first Norse settlements on the islands. Around 1,000 long-finned pilot whales are still killed annually, mainly during the summer. Today, dolphin meat is consumed in a small number of countries worldwide, which include Japan and Peru (where it is referred to as chancho marino, or "sea pork"). In some parts of the world, such as Taiji (in Japan) and the Faroe Islands, dolphins are traditionally considered food, and are killed in harpoon or drive hunts.
There have been human health concerns associated with the consumption of dolphin meat in Japan after tests showed that dolphin meat contained high levels of methylmercury. There are no known cases of mercury poisoning as a result of consuming dolphin meat, though the government continues to monitor people in areas where dolphin meat consumption is high. The Japanese government recommends that children and pregnant women avoid eating dolphin meat on a regular basis. Similar concerns exist with the consumption of dolphin meat in the Faroe Islands, where prenatal exposure to methylmercury and PCBs primarily from the consumption of pilot whale meat has resulted in neuropsychological deficits amongst children.
Ringed seals were once the main food staple for the Inuit. They are still an important food source for the people of Nunavut and are also hunted and eaten in Alaska. Seal meat is an important source of food for residents of small coastal communities. The seal blubber is used to make seal oil, which is marketed as a fish oil supplement. In 2001, two percent of Canada's raw seal oil was processed and sold in Canadian health stores.
In captivity
Aquariums
Cetaceans
Various species of dolphins are kept in captivity. These small cetaceans are more often than not kept in theme parks and dolphinariums, such as SeaWorld. Bottlenose dolphins are the most common species of dolphin kept in dolphinariums as they are relatively easy to train and have a long lifespan in captivity. Hundreds of bottlenose dolphins live in captivity across the world, though exact numbers are hard to determine. The dolphin "smile" makes them popular attractions, as this is a welcoming facial expression in humans; however, the smile is due to a lack of facial muscles and subsequent lack of facial expressions.
Organizations such as World Animal Protection and the Whale and Dolphin Conservation campaign against the practice of keeping cetaceans, particularly killer whales, in captivity. In captivity, they often develop pathologies, such as the dorsal fin collapse seen in 60–90% of male killer whales. Captives have vastly reduced life expectancies, on average only living into their twenties. In the wild, females who survive infancy live 46 years on average, and up to 70–80 years in rare cases. Wild males who survive infancy live 31 years on average, and up to 50–60 years. Captivity usually bears little resemblance to wild habitat, and captive whales' social groups are foreign to those found in the wild. Captive life is also stressful due to the requirement to perform circus tricks that are not part of wild killer whale behavior, as well as restricting pool size. Wild killer whales may travel up to in a day, and critics say the animals are too big and intelligent to be suitable for captivity. Captives occasionally act aggressively towards themselves, their tankmates, or humans, which critics say is a result of stress. Dolphins are often trained to do several anthropomorphic behaviors, including waving and kissing—behaviors wild dolphins would rarely do.
Pinnipeds
The large size and playfulness of pinnipeds make them popular attractions. Some exhibits have rocky backgrounds with artificial haul-out sites and a pool, while others have pens with small rocky, elevated shelters where the animals can dive into their pools. More elaborate exhibits contain deep pools that can be viewed underwater with rock-mimicking cement as haul-out areas. The most common pinniped species kept in captivity is the California sea lion as it is abundant and easy to train. These animals are used to perform tricks and entertain visitors. Other species popularly kept in captivity include the grey seal and harbor seal. Larger animals like walruses and Steller sea lions are much less common. Pinnipeds are popular attractions because they are "disneyfied", and consequently, people often anthropomorphize them with a curious, funny, or playful nature.
Some organizations, such as the Humane Society of the United States and World Animal Protection, object to keeping pinnipeds and other marine mammals in captivity. They state that the exhibits could not be large enough to house animals that have evolved to be migratory, and a pool could never replace the size and biodiversity of the ocean. They also oppose using sea lions for entertainment, claiming the tricks performed are "exaggerated variations of their natural behaviors" and distract the audience from the animal's unnatural environment.
Sea otter
Sea otters can do well in captivity, and are featured in over 40 public aquariums and zoos. The Seattle Aquarium became the first institution to raise sea otters from conception to adulthood with the birth of Tichuk in 1979, followed by three more pups in the early 1980s. In 2007, a YouTube video of two cute sea otters holding paws drew 1.5 million viewers in two weeks, and had over 20 million views . Filmed five years previously at the Vancouver Aquarium, it was YouTube's most popular animal video at the time, although it has since been surpassed. Otters are often viewed as having a "happy family life", but this is an anthropomorphism.
Sirenians
The oldest manatee in captivity was Snooty, at the South Florida Museum's Parker Manatee Aquarium in Bradenton, Florida. Born at the Miami Aquarium and Tackle Company on July 21, 1948, Snooty was one of the first recorded captive manatee births. He was raised entirely in captivity, and died at the age of 69. Manatees can also be viewed in a number of European zoos, such as the Tierpark in Berlin, the Nuremberg Zoo, in ZooParc de Beauval in France, and in the Aquarium of Genoa in Italy. The River Safari at Singapore features seven of them.
Military
Bottlenose dolphins and California sea lions are used in the United States Navy Marine Mammal Program (NMMP) to detect mines, protect ships from enemy soldiers, and recover objects. The Navy has never trained attack dolphins, as they would not be able to discern allied soldiers from enemy soldiers. There were five marine mammal teams, each purposed for one of the three tasks: MK4 (dolphins), MK5 (sea lions), MK6 (dolphins and sea lions), MK7 (dolphins) and MK8 (dolphins); MK is short for mark. The dolphin teams were trained to detect and mark mines either attached to the seafloor or floating in the water column, because dolphins can use their echolocative abilities to detect mines. The sea lion team retrieved test equipment such as fake mines or bombs dropped from planes usually out of reach of divers who would have to make multiple dives. MK6 protects harbors and ships from enemy divers, and was operational in the Gulf War and Vietnam War. The dolphins would swim up behind enemy divers and attach a buoy to their air tank, so that they would float to the surface and alert nearby Navy personnel. Sea lions would hand-cuff the enemy, and try to outmaneuver their counter-attacks.
The use of marine mammals by the Navy, even in accordance with the Navy's policy, continues to meet opposition. The Navy's policy says that only positive reinforcement is to be used while training the military dolphins, and that they be cared for in accordance with accepted standards in animal care. The inevitable stresses involved in training are topics of controversy, as their treatment is unlike the animals' natural lifestyle, especially towards their confined spaces when not training. There is also controversy over the use of muzzles and other inhibitors, which prevent the dolphins from foraging for food while working. The Navy states that this is to prevent them from ingesting harmful objects, but conservation activists say this is done to reinforce the trainers' control over the dolphins, who hand out food rewards. The means of transportation is also an issue for conservation activists, since they are hauled in dry carriers, and switching tanks and introducing the dolphin to new dolphins is potentially dangerous as they are territorial.
| Biology and health sciences | Mammals: General | Animals |
60258 | https://en.wikipedia.org/wiki/Sea%20lion | Sea lion | Sea lions are pinnipeds characterized by external ear flaps, long foreflippers, the ability to walk on all fours, short and thick hair, and a big chest and belly. Together with the fur seals, they make up the family Otariidae, eared seals. The sea lions have six extant and one extinct species (the Japanese sea lion) in five genera. Their range extends from the subarctic to tropical waters of the global ocean in both the Northern and Southern Hemispheres, with the notable exception of the northern Atlantic Ocean. They have an average lifespan of 20–30 years. A male California sea lion weighs on average about and is about long, while the female sea lion weighs and is long. The largest sea lions are Steller's sea lions, which can weigh and grow to a length of . Sea lions consume large quantities of food at a time and are known to eat about 5–8% of their body weight (about ) at a single feeding. Sea lions can move around in water and at their fastest they can reach a speed of about . Three species, the Australian sea lion, the Galápagos sea lion and the New Zealand sea lion, are listed as endangered.
Taxonomy
Sea lions are related to walruses and seals. Together with the fur seals, they constitute the family Otariidae, collectively known as eared seals. Until recently, sea lions were grouped under a single subfamily called Otariinae, whereas fur seals were grouped in the subfamily Arcocephalinae. This division was based on the most prominent common feature shared by the fur seals and absent in the sea lions, namely the dense underfur characteristic of the former. Recent genetic evidence, suggests Callorhinus, the genus of the northern fur seal, is more closely related to some sea lion species than to the other fur seal genus, Arctocephalus. Therefore, the fur seal/sea lion subfamily distinction has been eliminated from many taxonomies.
Nonetheless, all fur seals have certain features in common: the fur, generally smaller sizes, farther and longer foraging trips, smaller and more abundant prey items, and greater sexual dimorphism. All sea lions have certain features in common, in particular their coarse, short fur, greater bulk, and larger prey than fur seals. For these reasons, the distinction remains useful. The family Otariidae (Order Carnivora) contains the 15 extant species of fur seals and sea lions. Traditional classification of the family into the subfamilies Arctocephalinae (fur seals) and Otariinae (sea lions) is not supported, with the fur seal Callorhinus ursinus having a basal relationship relative to the rest of the family. This is consistent with the fossil record which suggests that this genus diverged from the line leading to the remaining fur seals and sea lions about 6 million years ago (mya). Similar genetic divergences between the sea lion clades as well as between the major Arctocephalus fur seal clades, suggest that these groups underwent periods of rapid radiation at about the time they diverged from each other. The phylogenetic relationships within the family and the genetic distances among some taxa highlight inconsistencies in the current taxonomic classification of the family.
Arctocephalus is characterized by ancestral character states such as dense underfur and the presence of double rooted cheek teeth and is thus thought to represent the most "primitive" line. It was from this basal line that both the sea lions and the remaining fur seal genus, Callorhinus, are thought to have diverged. The fossil record from the western coast of North America presents evidence for the divergence of Callorhinus about 6 mya, whereas fossils in both California and Japan suggest that sea lions did not diverge until years later.
Suborder Caniformia
Family Otariidae
Subfamily Arctocephalinae
Genus Arctocephalus (southern fur seal; eight species)
Genus Callorhinus (northern fur seal; one species)
Subfamily Otariinae
Genus Eumetopias
Steller's sea lion, E. jubatus
Genus Neophoca
Australian sea lion, N. cinerea
Genus Otaria
South American sea lion, O. flavescens
Genus Phocarctos
New Zealand sea lion or Hooker's sea lion, P. hookeri
Genus Zalophus
California sea lion, Z. californianus
Japanese sea lion, Z. japonicus – extinct (1950s)
Galapagos sea lion, Z. wollebaeki
Family Phocidae: true seals
Family Odobenidae: walrus
Physiology
Diving adaptations
There are many components that make up sea lion physiology and these processes control aspects of their behavior. Physiology dictates thermoregulation, osmoregulation, reproduction, metabolic rate, and many other aspects of sea lion ecology including but not limited to their ability to dive to great depths. The sea lions' bodies control heart rate, gas exchange, digestion rate, and blood flow to allow individuals to dive for a long period of time and prevent side effects of high pressure at depth.
The high pressures associated with deep dives cause gases such as nitrogen to build up in tissues which are then released upon surfacing, possibly causing death. One of the ways sea lions deal with the extreme pressures is by limiting the amount of gas exchange that occurs when diving. The sea lion allows the alveoli to be compressed by the increasing water pressure thus forcing the surface air into cartilage lined airway just before the gas exchange surface. This process prevents any further oxygen exchange to the blood for muscles, requiring all muscles to be loaded with enough oxygen to last the duration of the dive. However, this shunt reduces the amount of compressed gases from entering tissues therefore reducing the risk of decompression sickness. The collapse of alveoli does not allow for any oxygen storage in the lungs, however. This means that sea lions must mitigate oxygen use in order to extend their dives. Oxygen availability is prolonged by the physiological control of heart rate in sea lions. By reducing heart rate to well below surface rates, oxygen is saved by reducing gas exchange as well as reducing the energy required for a high heart rate. Bradycardia is a control mechanism to allow a switch from pulmonary oxygen to oxygen stored in the muscles which is needed when the sea lions are diving to depth. Another way sea lions mitigate the oxygen obtained at the surface in dives is to reduce digestion rate. Digestion requires metabolic activity and therefore energy and oxygen are consumed during this process; however, sea lions can limit digestion rate and decrease it by at least 54%. This reduction in digestion results in a proportional reduction in oxygen use in the stomach and therefore a correlated oxygen supply for diving. Digestion rate in these sea lions increases back to normal rates immediately upon resurfacing. Oxygen depletion limits dive duration, but carbon dioxide (CO2) build-up also plays a role in the dive capabilities of many marine mammals. After a sea lion returns from a long dive, CO2 is not expired as fast as oxygen is replenished in the blood, due to the unloading complications with CO2. However, having more than normal levels of CO2 in the blood does not seem to adversely affect dive behavior. Compared to terrestrial mammals, sea lions have a higher tolerance to storing CO2 which is what normally tells mammals that they need to breathe. This ability to ignore a response to CO2 is likely brought on by increased carotid bodies which are sensors for oxygen levels that let the animal know its available oxygen supply. Yet, the sea lions cannot avoid the effects of gradual CO2 build-up which eventually causes the sea lions to spend more time at the surface after multiple repeated dives to allow for enough built up CO2 to be expired.
Parasites and diseases
Galapagos sea lions (Zalophus wollebaeki) can be infected with Philophthalmus zalophi, an eye fluke. These infections have heavy impacts on the survival of juveniles. The disease appears to be compounded by global warming. The number of infectious stages of different parasites species has a strong correlation with temperature change, therefore it is essential to consider the correlation between the increasing number of parasitic infections and climate changes. The Galapagos Islands go through seasonal changes in sea surface temperatures, which consist of high temperatures from the beginning of January through the month of May and lower temperatures throughout the rest of the year. Parasites surfaced in large numbers when the sea temperature was at its highest. Furthermore, data was collected by capturing sea lions in order to measure and determine their growth rates. Their growth rates were noted along with the citings of parasites which were found under the eyelid. The shocking results were that sea lions are affected by the parasites from the early ages of 3 weeks old up until the age of 4 to 8 months. The parasites found in the eye fluke did serious damage to the eye. From the data collected, 21 of the 91 survived; with a total of 70 deaths in just a span of two years. The parasites are attacking the pups at such young ages and causing the pups to not reach the age of reproduction. The death rates of the pups is surpassing the fertility rate by far. Since most pups are unable to reach the age of reproduction, the population is not growing fast enough to keep the species out of endangerment. Other parasites, like Anisakis and heartworm, can also infect sea lions.
Australian sea lions (Neophoca cinerea) are also being affected by more frequent parasitic infections. The same method was used for the sea pups as on the Galapagos Islands, but in addition, the researchers in Australia took blood samples. The pups in Australia were being affected by hookworms, but they were also coming out in large numbers with warmer temperatures. New Zealand sea lion pups (Phocarctos hookeri) were also affected in really early ages by hookworms (Uncinaria). The difference is that in New Zealand researchers took the necessary steps and began treatment. The treatment seemed to be effective on the pups who have taken it. They found no traces of this infection afterwards. However, the percentage of pups who do have it is still relatively high at about 75%. Those pups who were treated had much better growth rates than those who did not. Overall parasites and hookworms are killing off enough pups to place them in endangerment. Parasites affect sea pups in various areas of the world. Reproductive success reduces immensely, survival methods, changes in health and growth have also been affected.
Similarly, climate change has resulted in increased toxic algae blooms in the oceans. These toxins are ingested by sardines and other fish which are then eaten by the sea lions, causing neurological damage and diseases such as epilepsy.
Gene expressions and diet
Gene expressions are being used more often to detect the physiological responses to nutrition, as well as other stressors. In a study done with four Steller sea lions (Eumetopias jubatus), three of the four sea lions underwent a 70-day trial which consisted of unrestricted food intake, acute nutritional stress, and chronic nutritional stress. Results showed that individuals under nutritional stress down-regulated some cellular processes within their immune response and oxidative stress. Nutritional stress was considered the most proximate cause of population decline in this species. In New Zealand sea lions, north-to south gradients driven by temperature differences were shown to be key factors in the prey mix. Adult California sea lions eat about 5% to 8% of their body weight per day ().
California sea lions feed mainly offshore in coastal areas. They eat a variety of prey—such as squid, anchovies, mackerel, rockfish, and sardines—found in upwelling areas. They also may take fish from commercial fishing gear, sport fishing lines, and fish passage facilities at dams and rivers.
Geographic variation
Geographic variation for sea lions have been determined by the observations of skulls of several Otariidae species; a general change in size corresponds with a change in latitude and primary productivity. Skulls of Australian sea lions from Western Australia were generally smaller in length whereas the largest skulls are from cool temperate localities. Otariidae are in the process of species divergence, much of which may be driven by local factors, particularly latitude and resources. Populations of a given species tend to be smaller in the tropics, increase in size with increasing latitude, and reach a maximum in sub-polar regions. In a cool climate and cold waters there should be a selective advantage in the relative reduction of body surface area resulting from increased size, since the metabolic rate is related more closely to body surface area than to body weight.
Breeding and population
Breeding methods and habits
Sea lions, with three groups of pinnipeds, have multiple breeding methods and habits over their families but they remain relatively universal. Otariids, or eared sea lions, raise their young, mate, and rest in more earthly land or ice habitats. Their abundance and haul-out behavior have a direct effect on their on land breeding activity. Their seasonal abundance trend correlates with their breeding period between the austral summer of January to March. Their rookeries populate with newborn pups as well as male and female otariids that remain to defend their territories. At the end of the breeding period males disseminate for food and rest while females remain for nurturing. Other points in the year consist of a mix of ages and genders in the rookeries with haul-out patterns varying monthly.
Steller sea lions, living an average of 15 to 20 years, begin their breeding season when adult males establish territories along the rookeries in early May. Male sea lions reach sexual maturity from ages 5 to 7 and do not become territorial until around 9 to 13 years of age. The females arrive in late May bringing in an increase of territorial defense through fighting and boundary displays. After a week births consist most usually of one pup with a perinatal period of 3 to 13 days.
Steller sea lions have exhibited multiple competitive strategies for reproductive success. Sea lion mating is often polygamous as males usually mate with different females to increase fitness and success, leaving some males to not find a mate at all. Polygamous males rarely provide parental care towards the pup. Strategies used to monopolize females include the resource-defense polygyny, or occupying important female resources. This involves occupying and defending a territory with resources or features attractive to females during sexually receptive periods. Some of these factors may include pupping habitat and access to water. Other techniques include potentially limiting access of other males to females.
Population
Otaria flavescens (South American sea lion) lives along the Chilean coast with a population estimate of 165,000. According to the most recent surveys in northern and southern Chile the sealing period of the middle twentieth century that left a significant decline in sea lion population is recovering. The recovery is associated with less hunting, otariids rapid population growth, legislation on nature reserves, and new food resources. Haul-out patterns change the abundance of sea lions at particular times of the day, month, and year. Patterns in migration relate to temperature, solar radiation, and prey and water resources. Studies of South American sea lions and other otariids document maximum population on land during early afternoon, potentially due to haul-out during high air temperatures. Adult and subadult males do not show clear annual patterns, maximum abundance being found from October to January. Females and their pups hauled-out during austral winter months of June to September.
Interactions with humans
South American sea lions have been greatly impacted by human exploitation. During the late Holocene period to the middle of the twentieth century, hunter-gatherers along the Beagle Channel and northern Patagonia had greatly reduced the number of sea lions due to their hunting of the species and exploitation of the species' environment. Although sealing has been put to a halt, in many countries, such as Uruguay, the sea lion population continues to decline because of the drastic effects humans have on their ecosystems. As a result, South American sea lions have been foraging at higher tropical latitudes than they did prior to human exploitation. Fishermen play a key role in the endangerment of sea lions. Sea lions rely on fish, like pollock, as a food source and have to compete with fishermen for it. When fishermen are successful at their job, they greatly reduce the sea lion's food source, which in turn endangers the species. Also, human presence and human recreational activities can cause sea lions to engage in violent and aggressive actions. When humans come closer than 15 meters of a sea lion, the sea lions' vigilance increases because of the disturbance of humans. These disturbances can potentially cause sea lions to have psychological stress responses that cause the sea lions to retreat, sometimes even abandon their locations, and decreases the amount of time sea lions spend hauling out.
New Zealand sea lions were also exploited from hunting and sealing, and as a result were extirpated from New Zealand's mainland for over 150 years, with their population being restricted to the subantarctic. In 1993, a female New Zealand sea lion gave birth on the mainland for the first time, and since then, they have slowly been recolonizing. These sea lions are the only pinnipeds that regularly move up to inland into forests. As consequence, they have been hit by cars on roads, deliberately killed, and been disturbed by dogs. Females need to move inland as a way to protect their pups, so roads, fences, residential areas, and private lands can inhibit their dispersal and breeding success. They have also adapted to commercial pine forests, and have given birth or nursed pups in residents' backyards and on golf courses. As one of the world's rarest sea lions, and an endangered and endemic species, efforts are being made to facilitate coexistence between them and humans.
Sea lion attacks on humans are rare, but when humans come within approximately , it can be very unsafe. In a highly unusual attack in 2007 in Western Australia, a sea lion leapt from the water and seriously mauled a 13-year-old girl surfing behind a speedboat. The sea lion appeared to be preparing for a second attack when the girl was rescued. An Australian marine biologist suggested that the sea lion may have viewed the girl "like a rag doll toy" to be played with. In San Francisco, where an increasingly large population of California sea lions crowds docks along San Francisco Bay, incidents have been reported in recent years of swimmers being bitten on the legs by large, aggressive males, possibly as territorial acts. In April 2015, a sea lion attacked a 62-year-old man who was boating with his wife in San Diego. The attack left the man with a punctured bone. In May 2017, a sea lion grabbed and pulled a girl into the water by her dress before retreating. The child was sitting on a pier side in British Columbia while tourists were illegally feeding the sea lions when the incident took place. She was pulled out of the water with minor injuries and received antibiotic prophylactic treatment for seal finger infection from the superficial bite injury.
There have also been documented events of sea lions assisting humans. One such notable instance of this is when Kevin Hines jumped off the Golden Gate Bridge in a suicide attempt and was helped to stay afloat by a sea lion until he was rescued by the Coast Guard.
Sea lions have also been a focus of tourism in Australia and New Zealand. One of the main sites to view sea lions is in the Carnac Island Nature Reserve near Perth in Western Australia. This tourist site receives over 100,000 visitors, many of whom are recreational boaters and tourists, who can watch the male sea lions haul out on to the shore. They have sometimes been called "the unofficial welcoming committee of the Galápagos Islands".
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| Biology and health sciences | Pinnipeds | Animals |
60261 | https://en.wikipedia.org/wiki/Pinniped | Pinniped | Pinnipeds (pronounced ), commonly known as seals, are a widely distributed and diverse clade of carnivorous, fin-footed, semiaquatic, mostly marine mammals. They comprise the extant families Odobenidae (whose only living member is the walrus), Otariidae (the eared seals: sea lions and fur seals), and Phocidae (the earless seals, or true seals), with 34 extant species and more than 50 extinct species described from fossils. While seals were historically thought to have descended from two ancestral lines, molecular evidence supports them as a monophyletic group (descended from one ancestor). Pinnipeds belong to the suborder Caniformia of the order Carnivora; their closest living relatives are musteloids (weasels, raccoons, skunks and red pandas), having diverged about 50 million years ago.
Seals range in size from the and Baikal seal to the and southern elephant seal. Several species exhibit sexual dimorphism. They have streamlined bodies and four limbs that are modified into flippers. Though not as fast in the water as dolphins, seals are more flexible and agile. Otariids primarily use their front limbs to propel themselves through the water, while phocids and walruses primarily use their hind limbs for this purpose. Otariids and walruses have hind limbs that can be pulled under the body and used as legs on land. By comparison, terrestrial locomotion by phocids is more cumbersome. Otariids have visible external ears, while phocids and walruses lack these. Pinnipeds have well-developed senses—their eyesight and hearing are adapted for both air and water, and they have an advanced tactile system in their whiskers or vibrissae. Some species are well adapted for diving to great depths. They have a layer of fat, or blubber, under the skin to keep warm in cold water, and, other than the walrus, all species are covered in fur.
Although pinnipeds are widespread, most species prefer the colder waters of the Northern and Southern Hemispheres. They spend most of their lives in water, but come ashore to mate, give birth, molt or to avoid ocean predators, such as sharks and orcas. Seals mainly live in marine environments but can also be found in fresh water. They feed largely on fish and marine invertebrates; a few, such as the leopard seal, feed on large vertebrates, such as penguins and other seals. Walruses are specialized for feeding on bottom-dwelling mollusks. Male pinnipeds typically mate with more than one female (polygyny), though the degree of polygyny varies with the species. The males of land-breeding species tend to mate with a greater number of females than those of ice breeding species. Male pinniped strategies for reproductive success vary between defending females, defending territories that attract females and performing ritual displays or lek mating. Pups are typically born in the spring and summer months and females bear almost all the responsibility for raising them. Mothers of some species fast and nurse their young for a relatively short period of time while others take foraging trips at sea between nursing bouts. Walruses are known to nurse their young while at sea. Seals produce a number of vocalizations, notably the barks of California sea lions, the gong-like calls of walruses and the complex songs of Weddell seals.
The meat, blubber and skin of pinnipeds have traditionally been used by indigenous peoples of the Arctic. Seals have been depicted in various cultures worldwide. They are commonly kept in captivity and are even sometimes trained to perform tricks and tasks. Once relentlessly hunted by commercial industries for their products, seals are now protected by international law. The Japanese sea lion and the Caribbean monk seal have become extinct in the past century, while the Mediterranean monk seal and Hawaiian monk seal are ranked as endangered by the International Union for Conservation of Nature. Besides hunting, pinnipeds also face threats from accidental trapping, marine pollution, climate change and conflicts with local people.
Etymology
The name "pinniped" derives from the Latin words and . The common name "seal" originates from the Old English word , which is in turn derived from the Proto-Germanic .
Taxonomy
The German naturalist Johann Karl Wilhelm Illiger was the first to recognize the pinnipeds as a distinct taxonomic unit; in 1811 he gave the name Pinnipedia to both a family and an order. American zoologist Joel Asaph Allen reviewed the world's pinnipeds in an 1880 monograph, History of North American pinnipeds, a monograph of the walruses, sea-lions, sea-bears and seals of North America. In this publication, he traced the history of names, gave keys to families and genera, described North American species and provided synopses of species in other parts of the world. In 1989, Annalisa Berta and colleagues proposed the unranked clade Pinnipedimorpha to contain the fossil genus Enaliarctos and modern seals as a sister group. Pinnipeds belong to the order Carnivora and the suborder Caniformia (known as dog-like carnivorans). Of the three extant families, the Otariidae and Odobenidae are grouped in the superfamily Otarioidea, while the Phocidae belong to the superfamily Phocoidea. There are 34 extant species of pinnipeds, and more than 50 fossil species of pinnipedimorphs.
Otariids are also known as eared seals due to their pinnae. These animals swim mainly using their well-developed fore-flippers. They can also "walk" on land by shifting their hind-flippers forward under the body. The front end of an otariid's frontal bone protrudes between the nasal bones, with a large and flattened supraorbital foramen. An extra spine splits the supraspinatous fossa and bronchi that are divided in the front. Otariids consist of two types: sea lions and fur seals; the latter typically being smaller, with pointier snouts, longer fore-flippers and heavier fur coats. Five genera and seven species (one now extinct) of sea lion are known to exist, while two genera and nine species of fur seal exist. While sea lions and fur seals have historically been considered separate subfamilies (Otariinae and Arctocephalinae respectively), genetic and molecular evidence has refuted this, indicating that the northern fur seal is basal to other otariids and the Australian sea lion and New Zealand sea lion are more closely related to Arctocephalus than to other sea lions.
Odobenidae has only one living member: the walrus. This animal is noticeable from its larger size (exceeded only by the elephant seals), nearly hairless skin, flattened snout and long upper canines, known as tusks. Like otariids, walruses can walk on land with their hind limbs. When moving in water, the walrus relies on its hind limbs for locomotion, while its forelimbs are used for steering. Also, it has no outer ears. The epipterygoid of the jaw is well developed and the back of the nasal bones are horizontal. In the feet, the calcaneuses protrude in the middle.
Phocids are known as true or "earless" seals. These animals lack outer ears and cannot position their hind-flippers to move on land, making them more cumbersome. This is because of their massive ankle bones and flatter heels. In water, true seals rely on the side-to-side motion of their hind-flippers and lower body to move forward. The phocid's skull has thickened mastoids, puffed up entotympanic bones, nasal bones with a pointed tip in the back and a non-existent supraorbital foramen. The hip has a more converse ilium. A 2006 molecular study supports the division of phocids into two monophyletic subfamilies: Monachinae, which consists of elephant seals, monk seals and Antarctic seals; and Phocinae, which consists of all the rest.
Evolution
One popular hypothesis suggested that pinnipeds are diphyletic (descended from two ancestral lines), with walruses and otariids sharing a recent common ancestor with bears; and phocids sharing one with Musteloidea. However, morphological and molecular evidence support a monophyletic origin. A 2021 genetic study found that pinnipeds are more closely related to musteloids. Pinnipeds split from other caniforms 50 million years ago (mya) during the Eocene. The earliest fossils of pinnipeds date back to the Late Oligocene. Fossil animals representing basal lineages include Puijila, of the Early Miocene in Arctic Canada. It resembled a modern otter, but shows evidence of quadrupedal swimming—retaining a form of aquatic locomotion that led to those employed by modern pinnipeds. Potamotherium, which lived in the same period in Europe, was similar to Puijila but more aquatic. The braincase of Potamotherium shows evidence that it used its whiskers to hunt, like modern seals. Both Puijila and Potamotherium fossils have been found in lake deposits, suggesting that seal ancestors were originally adapted for fresh water.
Enaliarctos, a fossil species of late Oligocene/early Miocene (24–22 mya) California, closely resembled modern pinnipeds; it was adapted to an aquatic life with flippers and a flexible spine. Its teeth were more like land predators in that they were more adapted for shearing. Its hind-flippers may have allowed it to walk on land, and it probably did not leave coastal areas as much as its modern relatives. Enaliarctos was likely more of a fore-flipper swimmer, but could probably swim with either pair. One species, Enaliarctos emlongi, exhibited notable sexual dimorphism, suggesting that this physical characteristic may have been an important driver of pinniped evolution. A closer relative of extant pinnipeds was Pteronarctos, which lived in Oregon 19–15 mya. As in modern seals, the maxilla or upper jaw bone of Pteroarctos intersects with the orbital wall. The extinct family Desmatophocidae lived 23–10 mya in the North Pacific. They had long skulls that with large orbits, interlocked zygomatic bones and rounded molars and premolars. They also were sexually dimorphic and may have been capable of swimming with both or either pair of flippers. They are grouped with modern pinnipeds, but there is debate as to whether they are more closely related to phocids or to otariids and walruses.
The ancestors of the Otarioidea and Phocidea diverged around 25 mya. Phocids are known to have existed for at least 15 million years, and molecular evidence supports a divergence of the Monachinae and Phocinae lineages around this time. The fossil genera Monotherium and Leptophoca of southeastern North America represent the earliest members of Monachinae and Phocinae respectively. Both lineages may have originated in the North Atlantic, and likely reached the Pacific via the Central American Seaway. Phocines mainly stayed in the Northern Hemisphere, while the monachines diversified southward. The lineages of Otariidae and Odobenidae split around 20 mya. The earliest fossil records of otariids are in North Pacific and dated to around 11 mya. Early fossil genera include Pithanotaria and Thalassoleon. The Callorhinus lineage split the earliest, followed by the Eumetopias/Zalophus lineage and then the rest, which colonized the Southern Hemisphere. The earliest fossils of Odobenidae—Prototaria of Japan and Proneotherium of Oregon—date to 18–16 mya. These primitive walruses had normal sized canines and fed on fish instead of mollusks. Later taxa like Gomphotaria, Pontolis and Dusignathus had longer canines on both the upper and lower jaw. The familiar long upper tusks developed in the genera Valenictus and Odobenus. The lineage of the modern walrus may have spread from the North Pacific to the North Atlantic through the Caribbean and Central American Seaway 8–5 mya, and then back to the North Pacific via the Arctic 1 mya, or to the Arctic and subsequently the North Atlantic during the Pleistocene.
Anatomy and physiology
Pinnipeds have streamlined, spindle-shaped bodies with small or non-existent ear flaps, rounded heads, short muzzles, flexible necks, limbs modified into flippers and small tails. The mammary glands and genitals can withdraw into the body. Seals are unique among carnivorans in that their orbital walls are mostly shaped by the maxilla and are not contained by certain facial bones. Compared to land carnivores, pinnipeds have fewer teeth, which are pointed and cone-shaped. They are adapted for holding onto slippery prey rather than shearing meat like the carnassials of other carnivorans. The walrus has unique tusks which are long upper canines.
Pinnipeds range in size from the and Baikal seal to the and southern elephant seal. Overall, they tend to be larger than other carnivores. Several species have male-biased sexual dimorphism that depends on how polygynous a species is: highly polygynous species like elephant seals are extremely sexually dimorphic, while less polygynous species have males and females that are closer in size, or, in the case of Antarctic seals, females are moderately bigger. Males of sexually dimorphic species also tend to have secondary sex characteristics, such as larger or more prominent heads, necks, chests, crests, noses/proboscises and canine teeth as well as thicker fur and manes. Though more polygynous species tend to be sexually dimorphic, some evidence suggests that size differences between the sexes originated due to ecological differences, with polygyny developing later.
Almost all pinnipeds have fur coats, the exception being the walrus, which is only sparsely covered. Even some fully furred species (particularly sea lions) are less furry than land mammals. Fur seals have lush coats consisting of an undercoat and guard hairs. In species that live on ice, young pups have thicker coats than adults. The individual hairs on the coat, known collectively as lanugo, can trap heat from sunlight and keep the pup warm. Pinnipeds are typically countershaded, and are darker colored dorsally and lighter colored ventrally, which serves to counter the effects of self-shadowing caused by light shining over the ocean water. The pure white fur of harp seal pups conceals them in their Arctic environment. Several species have clashing patterns of light and dark pigmentation. All fully furred species molt; the process of which may be quick or gradual depending on the species. Seals have a layer of subcutaneous fat, known as blubber, that is particularly thick in phocids and walruses. Blubber serves both to keep the animals warm and to provide energy and nourishment when they are fasting. It can constitute as much as 50% of a pinniped's mass. Newborn pups have a thin layer of blubber, but some species compensate for this with thick lanugos.
The simple stomach of pinnipeds is typical of carnivores. Most species have neither a cecum nor a clear demarcation between the small and large intestines; the large intestine is comparatively short and only slightly wider than the small intestine. Small intestine lengths range from 8 times (California sea lion) to 25 times (elephant seal) the body length. The length of the intestine may be an adaptation to frequent deep diving, allowing for more room in the digestive tract for partially digested food. An appendix is absent in seals. As in most marine mammals, the kidneys are divided into lobes and filter out excess salt.
Locomotion
Pinnipeds have two pairs of flippers on the front and back, the fore-flippers and hind-flippers. Their elbows and ankles are not externally visible. Pinnipeds are not as fast as cetaceans, typically swimming at compared to around for several species of dolphin. Seals are more agile and flexible, and some otariids, such as the California sea lion, can make dorsal turns as the back of their heads can touch their hind flippers. Pinnipeds have several adaptions for reducing drag. In addition to their streamlined bodies, they have smooth networks of muscle bundles in their skin that may increase laminar flow and cut through the water. The hair erector muscles are absent, so their fur can be streamlined as they swim.
When swimming, otariids rely on their fore-flippers for locomotion in a wing-like manner similar to penguins and sea turtles. Fore-flipper movement is not continuous, and the animal glides between each stroke. Compared to terrestrial carnivorans, the fore-limb bones of otariids are reduced in length, giving them less resistance at the elbow joint as the flippers flap; the hind-flippers maneuver them. Phocids and walruses swim by moving their hind-flippers and lower body from side to side, while their fore-flippers are mainly used for maneuvering. Some species leap out of the water, and "ride" waves.
Pinnipeds can move around on land, though not as well as terrestrial animals. Otariids and walruses are capable of turning their hind-flippers forward and under the body so they can "walk" on all fours. The fore-flippers move along a transverse plane, rather than the sagittal plane like the limbs of land mammals. Otariids create momentum by laterally swaying their heads and necks. Sea lions have been recorded climbing up flights of stairs. Phocids lack the ability to walk on their hind-flippers, and must flop and wriggle their bodies forward as their fore-flippers keep them stable. In some species, the fore-flippers may act like oars pushing against the ground. Phocids can move faster on ice, as they are able to slide.
Senses
The eyes of pinnipeds are relatively large for their size and are positioned near the front of the head. Only the smaller eyes of the walruses are located on each side of the head; since they forage at the bottom for sedentary mollusks. A seal's eye is suited for seeing both underwater and in air. Most of retina is equidistant around the spherical lens. The cornea has a flattened center where refraction does not change between air and water. The vascular iris has a strong dilator muscle. A contracted pupil is typically pear-shaped, although the bearded seal's is more horizontal. Compared to deep-diving elephant seals, the iris of shallower species, such as harbor seals and California sea lions, does not change much in size between contraction and expansion. Seals are able to see in darkness with a tapetum lucidum, a reflecting layer that increases sensitivity by reflecting light back through the rods.
On land, pinnipeds are near-sighted in dim light. This is reduced in bright light as the retracted pupil decreases the ability of the lens and the cornea to refract (bend) light. Polar living seals like the harp seal have corneas that are adapted to the bright light that reflects off snow and ice. As such, they do not suffer snow blindness. Pinnipeds appear to have limited color vision as they lack S-cones. Flexible eye movement has been documented in seals. The walrus can project its eyes out from its sockets in both a forward and upward direction due to its advanced extraocular muscles and absence of an orbital roof. The seal eye is durable as the corneal epithelium is hardened by keratin, and the sclera is thick enough to withstand the pressures of diving. Seals also secrete mucus from the lacrimal gland to protect their eyes. As in many mammals and birds, pinnipeds possess nictitating membranes.
The pinniped ear is adapted for hearing underwater, where it can hear sound frequencies of up to 70,000 Hz. In air, hearing is somewhat reduced in pinnipeds compared to many terrestrial mammals. While their airborne hearing sensitivity is generally weaker than humans', they still have a wide frequency range. One study of three species—the harbor seal, California sea lion and northern elephant seal—found that the sea lion was best adapted for airborne hearing, the elephant seal for underwater hearing and the harbor seal was equally adapted for both. Although pinnipeds have a fairly good sense of smell on land, it is useless under water as their nostrils are closed.
Pinnipeds have well-developed tactile senses. Compared to terrestrial mammals, the moustache-like whiskers or vibrissae of pinnipeds have ten times more nerve connections, allowing them to effectively detect vibrations in the water. These vibrations are generated, for example, when a fish swims through water. Detecting vibrations is useful when the animals are foraging, and may add to or even replace vision, particularly in darkness. Harbor seals can follow hydrodynamic paths made by other animals minutes earlier, similar to a dog following a scent trail, and even to discriminate the size and type of object responsible for the trail.
Unlike terrestrial mammals, such as rodents, pinnipeds do not sweep their whiskers over an object when examining it, but can protract the hairs forward while holding them steady, maximizing their detection. The vibrissa's angle relative to the flow seems to be the most important contributor to detection ability. The whiskers of some otariids grow quite long—those of the Antarctic fur seal can reach . Walruses have the most vibrissae, at 600–700 individual hairs. These are important when searching for prey along the bottom. In addition to foraging, whiskers may also play a role in navigation; spotted seals appear to use them to detect breathing holes in the ice.
Diving adaptations
To dive, a pinniped must first exhale much of the air out of its lungs and shut its nostrils and throat cartilages to protect the trachea. The airways are supported by cartilaginous rings and smooth muscle, and the chest muscles and alveoli can completely deflate during deeper dives. While land mammals generally cannot empty their lungs, pinnipeds can reinflate their lungs even after respiratory collapse. The middle ear contains sinuses that probably fill with blood during dives, preventing middle ear squeeze. The heart of a seal is moderately flattened to allow the lungs to deflate. The trachea is flexible enough to collapse under pressure. During deep dives, any remaining air in their bodies is stored in the bronchioles and trachea, which prevents them from experiencing decompression sickness, oxygen toxicity and nitrogen narcosis. In addition, seals can tolerate large amounts of lactic acid, which reduces skeletal muscle fatigue during intense physical activity.
The circulatory system of pinnipeds is large and elaborate; retia mirabilia line the inside of the trunk and limbs, allowing for greater oxygen storage during diving. As with other diving mammals, pinnipeds have large amounts of hemoglobin and myoglobin stored in their blood and muscles respectively. This allows them to stay submerged for long periods of time while still having enough oxygen. Deep-diving species such as elephant seals have blood volumes that represent up to 20% of their body weight. When diving, they reduce their heart rate, and blood flow is mostly restricted to the heart, brain and lungs. To keep their blood pressure stable, phocids have an elastic aorta that dissipates some of the energy of each heartbeat.
Thermoregulation
Pinnipeds keep warm by having large, thick bodies, insulating blubber and fur, and fast metabolism. Their idle body temperature is around against the ocean water. Metabolic rates of different species vary between 1.5 to 3 times that of land mammals. Also, the blood vessels in their flippers are adapted for countercurrent exchange; small veins surround arteries transporting blood from the body core, capturing heat from them. While blubber and fur keep the seal warm in water, they can also overheat the animal when it is on land. To counteract overheating, many species cool off by covering themselves in sand. Monk seals may even dig up the cooler layers. The northern fur seal cools off by panting.
Sleep
Pinnipeds spend many months at a time at sea, so they must sleep in the water. Scientists have recorded them sleeping for minutes at a time while slowly drifting downward in a belly-up orientation. Like other marine mammals, seals sleep in water with half of their brain awake so that they can detect and escape from predators, as well as surface for air without fully waking. When they are asleep on land, both sides of their brain go into sleep mode.
Distribution and habitat
Living pinnipeds are widespread in cold oceanic waters; particularly in the North Atlantic, the North Pacific and the Southern Ocean. By contrast, the consistently warm Indomalayan waters have no seals. Monk seals and some otariids live in tropical and subtropical waters. Seals usually require cool, nutrient-rich waters with temperatures lower than . Even in more tropical climates, lower temperatures and biological productivity may be provided by currents. Only monk seals live in waters that generally lack these features. The Caspian seal and Baikal seal are found in large landlocked bodies of water (the Caspian Sea and Lake Baikal respectively).
As a whole, pinnipeds can be found in a variety of aquatic habitats, mostly coastal water, but also open ocean, deep waters near offshore islands, brackish waters and even freshwater lakes and rivers. The Baikal seal is the only exclusively freshwater species. Pinnipeds also use a number of terrestrial habitats and substrates, both continental and island. In non-polar regions, they haul out on to rocky shores, sandy and pebble beaches, sandbanks, tidal flats or pools, and in sea caves. Some species also rest on man-made structures built along the coast or offshore. Pinnipeds may move further from the water using sand dunes or vegetation, or even rocky cliffs. New Zealand sea lions may travel to forests from the ocean. In polar regions, seals haul out on to both fast ice and drift ice. Some even den underneath the ice, particularly in pressure ridges and crevasses.
Behavior and life history
Pinnipeds have an amphibious lifestyle; they are mostly aquatic, but haul out to breed, molt, rest, sun or to avoid aquatic predators. Several species are known to migrate over vast distances, particularly in response to environmental changes. Elephant seals are at sea for most of the year and there are vast distances between their breeding and molting sites. The northern elephant seal is one of farthest mammalian migraters, traveling . Otariids tend to migrate less than phocids, especially tropical species. Traveling seals may reach their destination using geomagnetic fields, water and wind currents, solar and lunar positions and the temperature and chemical makeup of the water.
Pinnipeds may dive during foraging or to avoid predators. When foraging, for example, the Weddell seal typically dives for no more than 15 minutes and deep, but can dive for as long as 73 minutes and reach deep. Northern elephant seals often dive for as long as 20 minutes. They can also dive and for over an hour. The dives of otariids tend to be shorter and less deep. They typically last 5–7 minutes with average depths to . However, the New Zealand sea lion has been recorded diving to a maximum of and have submerged for as long as 12 minutes. The diet of walruses does not require them to dive very deep or very long. Pinnipeds generally live 25–30 years.
Foraging and predation
All pinnipeds are carnivorous and predatory. As a whole, they mostly feed on fish and cephalopods, but also consume crustaceans, bivalves, zooplankton and endothermic (warm-blooded) prey like sea birds. While most species have generalist diets, a few are specialists. Examples are krill-eating crabeater seals, crustacean-eating ringed seals, squid specialists like the Ross seal and southern elephant seal, and the bearded seal and walrus, which specialize on benthic invertebrates.
Pinnipeds may hunt solitarily or cooperatively. The former behavior is typical when hunting non-schooling fish, immobile or sluggish invertebrates and endothermic prey. Solitary foraging species usually hunt in coastal or shallow water. An exception to this is the northern elephant seal, which hunts deep in the open ocean for fish. In addition, walruses feed solitarily but are often near other walruses in small or large groups. For large schools of fish or squid, pinnipeds such as certain otariids hunt cooperatively in large groups, locating and herding their prey. Some species, such as California and South American sea lions, will hunt alongside sea birds and cetaceans.
Seals typically swallow their food whole, and will rip apart prey that is too big. The leopard seal, a prolific predator of penguins, is known to violently shake its prey to death. Complex serrations in the teeth of filter-feeding species, such as crabeater seals, allow water to leak out as they swallow their planktonic food. The walrus is unique in that it consumes its prey by suction feeding, using its tongue to suck the meat of a bivalve out of the shell. While pinnipeds mostly hunt in the water, South American sea lions are known to chase down penguins on land. Some species may swallow stones or pebbles for reasons not understood. Though they can drink seawater, pinnipeds get most of their fluid intake from their food.
Pinnipeds themselves are subject to predation. Most species are preyed on by the orca. To subdue and kill seals, orcas strike them with their heads or tails—the latter causing them to fly in the air—or simply bite into them and rip them apart. They are typically hunted by groups of 10 or fewer whales, but they are occasionally hunted by larger groups or by lone individuals. All age classes may be targeted, but pups most of all. Large sharks are another major predator of pinnipeds—usually the great white shark but also the tiger shark and mako shark. Sharks usually attack by ambushing them from below. Injured seals that escape are usually able to recover from their wounds. Otariids that have been targeted in the hindquarters are more likely to survive, while phocids are more likely to survive with forequarters injures. Pinnipeds are also preyed on by terrestrial and pagophilic predators. The polar bear is a major predator of Arctic seals and walruses, particularly pups. Bears may seek out seals, or simply wait for them to come by. Other terrestrial predators include cougars, brown hyenas and various species of canids, which mostly target the young.
Pinnipeds lessen the chance of predation by gathering in groups. Some species are capable of inflicting damaging wounds on their attackers with their sharp canines. Adult walruses are particularly risky prey for polar bears. When out at sea, northern elephant seals dive out of the reach of surface-hunting orcas and white sharks. In the Antarctic, which lacks terrestrial predators, pinniped species spend more time on the ice than their Arctic counterparts.
Interspecific predation among pinnipeds does occur. The leopard seal is known to prey on many other species, especially the crabeater seal. Leopard seals typically target crabeater pups, particularly from November to January. Older crabeater seals commonly bear scars from failed leopard seal attacks; a 1977 study found that 75% of a sample of 85 individual crabeaters had these scars. Walruses, despite being specialized for feeding on bottom-dwelling invertebrates, occasionally prey on Arctic seals. They kill their prey with their long tusks and eat their blubber and skin. Steller sea lions have been recorded eating harbor seals, northern fur seals and California sea lions, particularly pups and small adults. New Zealand sea lions feed on pups of some fur seal species, and the South American sea lion may prey on South American fur seals.
Reproductive behavior
The mating system of pinnipeds varies from extreme polygyny to serial monogamy. Species that breed on land are usually more polygynous, as females gather in large aggregations and males are able to mate with them as well as defend them from rivals. These species include elephant seals, grey seals and most otariids. Land-breeding pinnipeds prefer to mate on islands where there are fewer land predators. Suitable islands are in short supply and tend to be crowded. Since the land they breed on is fixed, females return to the same sites for many years. The males arrive earlier in the season and wait for them. The males stay on land to monopolize females; and may fast for months as they would lose their position if they went to feed at sea. Polygynous species also tend to be extremely sexual dimorphic in favor of males. This dimorphism manifests itself in larger chests and necks, longer canines and denser fur—all traits that equip males for combat. Larger males have more blubber and thus more energy reserves for fasting.
Other seals, like the walrus and most phocids, breed on ice and copulate in the water—a few land-breeding species also mate in water. Females of these species tend to be more spaced out and there is less site fidelity, since ice is less stable than solid land. Hence polygyny tends to be weaker in ice-breeding species. An exception to this is the walrus, whose distribution of food forces females closer together. Pinnipeds that breed on fast ice tend to cluster together more than those that breed on drift ice. Seals that breed on ice tend to have little or no sexual dimorphism. In Antarctic seals, there is some size bias in favor of females. Walruses and hooded seals are unique among ice-breeding species in that they have pronounced sexual dimorphism in favor of males.
Adult male pinnipeds have several strategies to ensure reproductive success. Otariid males gain access to females by establishing territories where females can bask and give birth and contain valuable resources such as shade, tide pools or access to water. Territories are usually marked by natural barriers, and some may be fully or partially underwater. Males defend their territorial boundaries with threatening vocalizations and postures, but physical fights are usually not very violent, and are mostly limited to early in the season. Individuals also return to the same territorial site each breeding season. In certain species, like the Steller sea lion and northern fur seal, a dominant male can maintain a territory for as long as 2–3 months. Females can usually move freely between territories and males are unable to coerce females who are intent on leaving, but in some species such as the northern fur seal, South American sea lion and Australian sea lion, males keep females in their territories with threatening displays and even violence. In some phocid species, like the harbor seal, Weddell seal and bearded seal, the males establish "maritories" and patrol and defend the waters bordering female haul-out areas, waiting for a female to enter. These are also maintained by vocalizations. The maritories of Weddell seal males include entries to female breathing holes in the ice.
Lek systems are known to exist among some populations of walruses. These males gather near female herds and try to attract them with elaborate courtship displays and vocalizations. Lekking may also exist among California sea lions, South American fur seals, New Zealand sea lions and harbor seals. In some species, including elephant seals, grey seals and non-lekking walruses, males will try to lay claim to the desired females and defend them from rivals. Elephant seal males, in particular, establish dominance hierarchies via displays and fights, with the highest ranking males having a near monopoly on reproductive success. An alpha male can have a harem of 100 females. Grey seal males usually place themselves among a cluster of females whose members may change over time, while males of some walrus populations guard female herds. Male ringed, crabeater, spotted and hooded seals follow and defend nearby females and mate with them when they reach estrus. These may be lone females or small groups. South American sea lions are considered to be both a territory-defending and female-defending species. Males start the season establishing and defending territories but then claim and defend females when they arrive.
Younger or subdominant male pinnipeds may attempt to achieve reproductive success in other ways including sneakiness, harassment of females or even coordinated disruption of the colony. Female pinnipeds do appear to have some choice in mates, particularly in lek-breeding species like the walrus, but also in elephant seals where the males try to dominate all the females that they want to mate with. When a female elephant seal or grey seal is mounted by an unwanted male, she tries to resist and get away. This commotion attracts other males to the scene, and the most dominant will take over and mate with female himself. Dominant female elephant seals stay in the center of the colony where they in the domain of a more dominant male, while marginal females are left with subordinates. Female Steller sea lions may solicit their territorial males for mating.
Birth and parenting
Except for the walrus, which has five- to six-year gaps between births, female pinnipeds enter estrus shortly after they give birth, and can thus produce pups every year. All species have delayed implantation, in which the embryo does not enter the uterus for weeks or months. Delayed implantation allows the female to wait until conditions are right for birthing. Gestation in seals (including delayed implantation) typically lasts a year. For most species, birthing takes place in spring and summer. Usually, single pups are born; twins are rare and have high mortality rates. Pups of most species are born relatively developed and precocial.
Pinniped milk has "little to no lactose". Mother pinnipeds have different strategies for maternal care and lactation. Phocids such as elephant seals, grey seals and hooded seals have a lactation period that lasts days or weeks, during which they fast and nurse their pups on land or ice. The milk of these species consists of up to 60% fat, allowing the young to grow quickly. Each day until they are weaned, northern elephant seal pups gain . Some pups gain weight more quickly than others by stealing extra milk from other mothers. Alloparenting occurs in these fasting species; while most northern elephant seal mothers nurse their own pups and reject nursings from alien pups, some do accept alien pups with their own.
For otariids and some phocids like the harbor seal, mothers fast and nurse their pups for a few days at a time. In between nursing bouts, the females forage at sea while the young stay behind onshore. If there is enough food close to shore, a female can be gone for as little as a day, but otherwise may be at sea for as long as three weeks. Lactation in otariids may last 6–11 months; in the Galápagos fur seal it can last up to three years. Pups of these species are weaned at heavier weights than their phocid counterparts, but the latter grow quicker. Walruses are unique in that mothers nurse their young at sea. Young pinnipeds typically start swimming on their own and some species can even swim as newborns. Young may wait days or weeks before entering the water; elephant seals start swimming weeks after weaning.
Male pinnipeds generally play little role in raising the young. Male walruses may help inexperienced young as they learn to swim, and have even been recorded caring for orphans. When a group is threatened, all the adults may protect the young. Male California sea lions have been observed to help shield swimming pups from predators. Males can also pose threats to the safety of pups, particularly during fights. Pups of some species may be abducted, assaulted and killed by males.
Communication
Pinnipeds can produce a number of vocalizations. While most vocals are audible to the human ear, Weddell seals have been recorded in Antarctica making ultrasonic calls underwater. In addition, the vocals of northern elephant seals may produce infrasonic vibrations. Vocals are produced both in air and underwater; the former are more common among otariids and the latter among phocids. Antarctic seals are noisier on land or ice than Arctic seals due to the absence of polar bears. Male vocals are usually deeper than those of the females. Vocalizations are particularly important during the breeding seasons. Dominant male elephant seals display their status and power with "clap-threats" and loud drum-like calls that may be modified by the proboscis. Male otariids have strong barks, growls and roars. Male walruses are known to produce gong-like calls when attempting to attract females, these are amplified underwater with inflatable throat sacs.
The Weddell seal has perhaps the most extensive vocal repertoire, producing both airborne and underwater sounds. Trilling, gluping, chirping, chugging and knocking are some examples of the calls produced underwater. When warning other seals, the calls may be pronounced by "prefixes" and "suffixes". The underwater vocals of Weddell seals can last 70 seconds, which is long for a marine mammal call. Some calls have about seven rhythm patterns and could be categorized as "songs". Similar calls have been recorded in other Antarctic seals and in bearded seals. In some pinniped species, there appear to be regional dialects or even individual variations in vocalizations. These differences are likely important for territorial males becoming accustomed to their neighbors (dear enemy effect) and mothers and pups who need to remain in contact on crowded beaches. Female seals emit a "pulsed, bawling" contact call, while pups respond by squawking. Contact calls are particularly important for otariid mothers returning from sea. Other vocalizations produced by seals include grunts, rasps, rattles, creaks, warbles, clicks and whistles.
Non-vocal communication is not as common in pinnipeds as in cetaceans. Nevertheless, when they feel threatened, hauled-out harbor seals and Baikal seals may slap themselves with their flippers to create a warning sound. Teeth chattering, hisses and exhalations are also made as aggressive warnings by pinnipeds. Visual displays also occur: Ross seals resting on the ice will show the stripes on their chests and bare their teeth to a perceived threat, while swimming Weddell seals will make an S-shaped posture to intimidate rivals under the ice. Male hooded seals use their inflatable nasal membranes to display to and attract females.
Intelligence
In a match-to-sample task study, a single California sea lion was able to demonstrate an understanding of symmetry, transitivity and equivalence; a second seal was unable to complete the tasks. They demonstrate the ability to understand simple syntax and commands when taught an artificial sign language, though they only rarely used the signs semantically or logically. In 2011, a captive California sea lion named Ronan was recorded bobbing its head in synchrony to musical rhythms. This "rhythmic entrainment" was previously seen only in humans, parrots and other birds possessing vocal mimicry. Adult male elephant seals can recognize each other's vocalizations by remembering the rhythm and timbre. In the 1970s, a captive harbor seal named Hoover was trained to imitate human speech and laughter.
For sea lions used in entertainment, trainers toss a ball at the animal or simply place the object on its nose, so it will eventually understand the behavior desired. A sea lion may need a year of training before it can publicly perform. Its long-term memory allows it to perform a trick after as much as three months of non-performance.
Human relations
In culture
Various human cultures have for millennia depicted pinnipeds. In Homer's Odyssey, the sea god Proteus shepherds a colony of seals. In northern Scotland, Celts of Orkney and the Hebrides believed in selkies—seals that could change into humans and walk on land. In Inuit mythology, they are associated with the goddess Sedna, who sometimes transformed into a seal. It was believed that marine mammals, including seals, came from her severed fingers. In modern culture, pinnipeds are thought of as cute, playful and comical figures.
In captivity
Pinnipeds can be found in facilities around the world, as their size and playfulness make them popular attractions. Seals have been kept in captivity since at least Ancient Rome and their trainability was noticed by Pliny the Elder. Zoologist Georges Cuvier noted during the 19th century that wild seals show considerable affection for humans and stated that they are second only to some monkeys among wild animals in their easy tamability. Francis Galton noted in his seminal work on domestication that seals were a spectacular example of an animal that would most likely never be domesticated, despite their friendliness, survivability and "desire for comfort", because they serve no practical use for humans.
Some modern exhibits have a pool with artificial haul-out sites and a rocky background, while others have seals housed in shelters located above a pool which they can jump into. More elaborate exhibits contain deep pools that can be viewed underwater with rock-mimicking cement as haul-out areas. The most popular captive pinniped is the California sea lion, due to its trainability and adaptability. Other commonly kept species include the grey seal and harbor seal. Larger animals like walruses and Steller sea lions are much less common. Some organizations, such as the Humane Society of the United States and World Animal Protection, object to keeping marine mammals in captivity. They state that the exhibits could not be large enough to house animals that have evolved to be migratory, and a pool could never replace the size and biodiversity of the ocean. They also state that the tricks performed for audiences are "exaggerated variations of their natural behaviors" and distract the people from the animal's unnatural environment.
California sea lions are used in military applications by the U.S. Navy Marine Mammal Program, including detecting naval mines and enemy divers. In the Persian Gulf, the animals have been trained to swim behind divers approaching a U.S. naval ship and attach a clamp with a rope to the diver's leg. Navy officials say that the sea lions can do this in seconds, before the enemy realizes what happened. Organizations like PETA believe that such operations put the animals in danger. The Navy insists that the sea lions are removed once their job is done.
Hunting
Humans have hunted seals since the Stone Age. Originally, seals were merely hit with clubs during haul-out. Eventually, more lethal weapons were used, like spears and harpoons. They were also trapped in nets. The use of firearms in seal hunting during the modern era drastically increased the number of killings. Pinnipeds are typically hunted for their meat and blubber. The skins of fur seals and phocids are made into coats, and the tusks of walruses have been used as ivory. There is a distinction made between the subsistence hunting of seals by indigenous peoples of the Arctic and commercial hunting: subsistence hunters depend on seal products for survival. National and international authorities have given special treatment to aboriginal hunters since their methods of killing are seen as more sustainable and smaller in scope. However indigenous people have recently used more modern technology and are profiting more from seal products in the marketplace. Some anthropologists argue that the term "subsistence" should also apply to these activities, as long as they are local in scale. More than 100,000 phocids (especially ringed seals) as well as around 10,000 walruses are harvested annually by native hunters.
Commercial sealing rivaled whaling as an important industry throughout history. Harvested species included harp seals, hooded seals, Caspian seals, elephant seals, walruses and all species of fur seal. After the 1960s, the harvesting of seals decreased substantially as an industry after the Canadian government implemented measures to protect female seals and restrict the hunting season. Several species that were commercially exploited have rebounded in numbers; for example, Antarctic fur seals may have reached their pre-harvesting numbers. The northern elephant seal nearly went extinct in the late 19th century, with only a small population remaining on Guadalupe Island. It has since recolonized much of its historic range, but has a population bottleneck. Conversely, the Mediterranean monk seal was extirpated from much of the Mediterranean and its current range is still limited.
Several species of pinniped continue to be exploited. The Convention for the Conservation of Antarctic Seals protects species within the Antarctic and surrounding waters, but allows restricted hunting of crabeater seals, leopard seals and Weddell seals. Weddell seal hunting is forbidden between September and February if the animal is older than a year, to ensure healthy population growth. The Government of Canada permits the hunting of harp seals. This has been met with controversy and debate. Proponents of seal hunts insist that the animals are killed humanely and the white-coated pups are not taken, while opponents argue that it is irresponsible to kill harp seals as they are already threatened by declining habitat.
The Caribbean monk seal has been killed and exploited by European settlers and their descendants since 1494, starting with Christopher Columbus himself. The seals were easy targets for organized sealers, fishermen, turtle hunters and buccaneers because they evolved with little pressure from terrestrial predators and were thus "genetically tame". In the Bahamas, as many as 100 seals were slaughtered in one night. The species was considered to be already extinct by the mid-nineteenth century until a small colony was found near the Yucatán Peninsula in 1866. Seal killings continued, and the last reliable report of the animal alive was in 1952 at Serranilla Bank. The IUCN declared it extinct in 1996. The Japanese sea lion was common around the Japanese islands, but overexploitation and competition from fisheries drastically decreased the population in the 1930s. The last recorded individual was a juvenile in 1974.
Conservation issues
As of 2021, the International Union for Conservation of Nature (IUCN) recognizes 36 pinniped species. With the Japanese sea lion and the Caribbean monk seal recently extinct, ten more are considered at risk. They are ranked as:
"Endangered": Hawaiian monk seal, Mediterranean monk seal, Galápagos fur seal, Australian sea lion, New Zealand sea lion, Caspian seal, and Galápagos sea lion.
"Vulnerable": northern fur seal, hooded seal, and walrus.
Pinnipeds face various threats. They are unintentionally caught in fishing nets by commercial fisheries and accidentally swallow fishing hooks. Gillnetting and Seine netting is a significant cause of mortality in seals and other marine mammals. Species commonly entangled include California sea lions, Hawaiian monk seals, northern fur seals and brown fur seals. Pinnipeds are also affected by marine pollution. Organic chemicals tend to accumulate in these animals since they are high in the food chain and have large reserves of blubber. Lactating mothers can pass the toxins on to their young. These pollutants can cause gastrointestinal cancers, and decreased fertility and immunity to infectious diseases. Other man-made threats include habitat destruction by oil and gas exploitation, encroachment by boats, and underwater noise.
Species that live in polar habitats are vulnerable to the effects of climate change on oceans, particularly declines in sea ice. In 2010 and 2011, sea ice in the Northwest Atlantic was at or near an all-time low and harp seals as well as ringed seals that bred on thin ice saw increased death rates. In the Antarctic, the decreased duration and extent of the sea ice and nutrient availability could potentially reduce the survival of Weddell seal pups and may have important implications for population growth rates. Antarctic fur seals in South Georgia in the South Atlantic saw major decreases over a 20-year study, during which scientists measured increased sea surface temperature anomalies.
Some species have become so numerous that they conflict with local people. In the United States, pinnipeds are protected under the Marine Mammal Protection Act of 1972 (MMPA). Since that year, California sea lion populations have risen to 250,000. These animals began exploiting more man-made environments, like docks, for haul-out sites. Many docks are not designed to withstand the weight of several resting sea lions. Wildlife managers have used various methods to control the animals, and some city officials have redesigned docks so they can better resist sea lion use. Inland-living New Zealand sea lions face unique human conflicts such as road mortality and run-ins with human infrastructure. Seals also conflict with fisherman. In 2007, MMPA was amended to permit the lethal removal of sea lions from salmon runs at Bonneville Dam. In the 1980s and 1990s, South African politicians and fishermen demanded that brown fur seals be culled, believing that the animals competed with commercial fisheries. Scientific studies found that culling fur seals would actually have a negative effect on the fishing industry, and the culling option was dropped in 1993.
| Biology and health sciences | Pinnipeds | null |
60264 | https://en.wikipedia.org/wiki/Mephitidae | Mephitidae | Mephitidae is a family of mammals comprising the skunks and stink badgers. They are noted for the great development of their anal scent glands, which they use to deter predators. Skunks were formerly classified as a subfamily of the Mustelidae (the weasel family); however, in the 1990s, genetic evidence caused skunks to be treated as a separate family. Similarly, the stink badgers had been classified with badgers, but genetic evidence shows they share a more recent common ancestor with skunks, so they are now included in the skunk family. A 2017 study using retroposon markers indicated that they are most closely related to the Ailuridae (red pandas and allies) and Procyonidae (raccoons and allies).
There are twelve extant species of mephitids in four genera: Conepatus (hog-nosed skunks, four species); Mephitis (the hooded and striped skunks, two species); Mydaus (stink badgers, two species); and Spilogale (spotted skunks, four species). The two stink badgers in the genus Mydaus inhabit Indonesia, Brunei, Malaysia and the Philippines; the other members of the family inhabit the Americas, ranging from Canada to central South America. All other mephitids are extinct, known through fossils, including those from Eurasia.
In taxonomic order, the living species of Mephitidae are:
Genera
| Biology and health sciences | Other carnivora | Animals |
60284 | https://en.wikipedia.org/wiki/Fountain%20pen | Fountain pen | A fountain pen is a writing instrument that uses a metal nib to apply water-based ink, or special pigment ink—suitable for fountain pens—to paper. It is distinguished from earlier dip pens by using an internal reservoir to hold ink, eliminating the need to repeatedly dip the pen in an inkwell during use. The pen draws ink from the reservoir through a feed to the nib and deposits the ink on paper via a combination of gravity and capillary action. Filling the reservoir with ink may be achieved manually, via the use of an eyedropper or syringe, or via an internal filling mechanism that creates suction (for example, through a piston mechanism) or a vacuum to transfer ink directly through the nib into the reservoir. Some pens employ removable reservoirs in the form of pre-filled ink cartridges.
History
Early prototypes of reservoir pens
According to Qadi al-Nu'man al-Tamimi () in his Kitab al-Majalis wa 'l-musayarat, the Fatimid caliph Al-Mu'izz li-Din Allah in Arab Egypt demanded a pen that would not stain his hands or clothes, and was provided with a pen that held ink in a reservoir, allowing it to be held upside-down without leaking.
There is compelling evidence that a working fountain pen was constructed and used during the Renaissance by artist and inventor Leonardo da Vinci. Leonardo's journals contain drawings with cross-sections of what appears to be a reservoir pen that works by both gravity and capillary action. Historians also took note of the fact that the handwriting in the inventor's surviving journals is of a consistent contrast throughout, rather than exhibiting the characteristic fading pattern typical of a quill pen caused by expending and re-dipping. While no physical item survives, several working models were reconstructed in 2011 by artist Amerigo Bombara that have since been put on display in museums dedicated to Leonardo.
European reservoir models
The fountain pen was available in Europe in the 17th century and is shown by contemporary references. In Deliciae Physico-Mathematicae (a 1636 magazine), German inventor Daniel Schwenter described a pen made from two quills. One quill served as a reservoir for ink inside the other quill. The ink was sealed inside the quill with cork. Ink was squeezed through a small hole to the writing point. In 1663 Samuel Pepys referred to a metal pen "to carry ink". Noted Maryland historian Hester Dorsey Richardson (1862–1933) documented a reference to "three silver fountain pens, worth 15 shillings" in England during the reign of Charles II, c. 1649–1685. By the early 18th century such pens were already commonly known as "fountain pens". Hester Dorsey Richardson also found a 1734 notation made by Robert Morris the elder in the ledger of the expenses of Robert Morris the younger, who was at the time in Philadelphia, for "one fountain pen". Perhaps the best-known reference, however, is that of Nicholas Bion (1652–1733), whose illustrated description of a "plume sans fin" was published in 1709 in his treatise published in English in 1723 as "The Construction and Principal Uses of Mathematical Instruments". The earliest datable pen of the form described by Bion is inscribed 1702, while other examples bear French hallmarks as late as the early 19th century.
First patents
Progress in developing a reliable pen was slow until the mid-19th century because of an imperfect understanding of the role that air pressure plays in the operation of pens. Furthermore, most inks were highly corrosive and full of sedimentary inclusions. The first English patent for a fountain pen was issued in May 1809 to Frederick Fölsch, with a patent covering (among other things) an improved fountain pen feed issued to Joseph Bramah in September 1809. John Scheffer's patent of 1819 was the first design to see commercial success, with a number of surviving examples of his "Penographic" known. Another noteworthy pioneer design was John Jacob Parker's, patented in 1832 – a self-filler with a screw-operated piston. The Romanian inventor Petrache Poenaru received a French patent on May 25, 1827, for the invention of a fountain pen with a barrel made from a large swan quill.
Mass-manufactured nibs
In 1828, Josiah Mason improved a cheap and efficient slip-in nib in Birmingham, England, which could be added to a fountain pen and in 1830, with the invention of a new machine, William Joseph Gillott, William Mitchell, and James Stephen Perry devised a way to mass manufacture robust, cheap steel pen nibs (Perry & Co.). This boosted the Birmingham pen trade and by the 1850s, more than half the steel-nib pens manufactured in the world were made in Birmingham. Thousands of skilled craftsmen were employed in the industry. Many new manufacturing techniques were perfected, enabling the city's factories to mass-produce their pens cheaply and efficiently. These were sold worldwide to many who previously could not afford to write, thus encouraging the development of education and literacy.
New patents and inventions
In 1848, American inventor Azel Storrs Lyman patented a pen with "a combined holder and nib". In 1849 Scottish inventor Robert William Thomson invented the refillable fountain pen. From the 1850s, there was a steadily accelerating stream of fountain pen patents and pens in production. However, it was only after three key inventions were in place that the fountain pen became a widely popular writing instrument. Those were the iridium-tipped gold nib, hard rubber, and free-flowing ink.
The first fountain pens making use of all these key ingredients appeared in the 1850s. In the 1870s Duncan MacKinnon, a Canadian living in New York City, and Alonzo T. Cross of Providence, Rhode Island, created stylographic pens with a hollow, tubular nib and a wire acting as a valve. Stylographic pens are now used mostly for drafting and technical drawing but were very popular in the decade beginning in 1875. In the 1880s the era of the mass-produced fountain pen finally began. The dominant American producers in this pioneer era were Waterman, of New York City, and Wirt, based in Bloomsburg, Pennsylvania. Waterman soon outstripped Wirt, along with many companies that sprang up to fill the new and growing fountain pen market. Waterman remained the market leader until the early 1920s.
At this time, fountain pens were almost all filled by unscrewing a portion of the hollow barrel or holder and inserting the ink by means of a dropper – a slow and messy procedure. Pens also tended to leak inside their caps and at the joint where the barrel opened for filling. Now that the materials' problems had been overcome and the flow of ink while writing had been regulated, the next problems to be solved were the creation of a simple, convenient self-filler and the problem of leakage. Self-fillers began to gain in popularity around the turn of the century; the most successful of these was probably the Conklin crescent-filler, followed by A. A. Waterman's twist-filler. The tipping point, however, was the runaway success of Walter A. Sheaffer's lever-filler, introduced in 1912, paralleled by Parker's roughly contemporary button-filler.
Pen leakage
Meanwhile, many inventors turned their attention to the problem of leakage. Some of the earliest solutions to this problem came in the form of a "safety" pen with a retractable point that allowed the ink reservoir to be corked like a bottle. Horton, Moore, and Caw's were the earliest makers of such pens, all starting in the 1890s.
In 1898, George Safford Parker released the Parker Jointless, so named because its barrel was single-piece with no section joint to leak and stain the writer's fingers. The nib and feed assembly fit into the barrel's end like a cork stopper.
In 1908, Waterman began marketing a popular safety pen of its own. For pens with non-retractable nibs, the adoption of screw-on caps with inner caps that sealed around the nib by bearing against the front of the section effectively solved the leakage problem (such pens were also marketed as "safety pens", as with the Parker Jack-Knife Safety and the Swan Safety Screw-Cap).
Further innovation
In Europe, the German office supplies company Gunther Wagner, founded in 1838, introduced their Pelikan in 1929, the first modern screw piston-filling fountain pen. This was based upon the acquisition of patents for solid-ink fountain pens from the factory of Slavoljub Penkala from Croatia (patented 1907, in mass production since 1911), and the patent of the Hungarian Theodor Kovacs for the modern piston filler by 1925.
The decades that followed saw many technological innovations in the manufacture of fountain pens. Celluloid gradually replaced hard rubber, which enabled production in a much wider range of colors and designs. At the same time, manufacturers experimented with new filling systems. The inter-war period saw the introduction of some of the most notable models, such as the Parker Duofold and Vacumatic, Sheaffer's Lifetime Balance series, and the Pelikan 100.
During the 1940s and 1950s, fountain pens retained their dominance: early ballpoint pens were expensive, were prone to leaks and had irregular inkflow, while the fountain pen continued to benefit from the combination of mass production and craftsmanship. (Bíró's patent, and other early patents on ball-point pens often used the term "ball-point fountain pen," because at the time the ball-point pen was considered a type of fountain pen; that is, a pen that held ink in an enclosed reservoir.)
This period saw the launch of innovative models such as the Parker 51, the Aurora 88, the Sheaffer Snorkel, and the Eversharp Skyline, while the Esterbrook J series of lever-fill models with interchangeable steel nibs offered inexpensive reliability to the masses.
Popular usage
By the 1960s, refinements in ballpoint pen production gradually ensured its dominance over the fountain pen for casual use. Although cartridge-filler fountain pens are still in common use in France, Italy, Germany, Austria, India, and the United Kingdom, and are widely used by young students in most private schools in England, at least one private school in Scotland, and public elementary schools in Germany, a few modern manufacturers (especially Conway Stewart, Montblanc, Graf von Faber-Castell, and Visconti) now depict the fountain pen as a collectible item or a status symbol, rather than an everyday writing tool. However, fountain pens continue to have a growing following among many who view them as superior writing instruments due to their relative smoothness and versatility. Retailers continue to sell fountain pens and inks for casual and calligraphic use. Recently, fountain pens have made a resurgence, with some retailers, such as Goulet Pens, saying it is because of renewed consumer interest in analog products. This has led to a new wave of casual use fountain pens and custom ink manufacturers, who utilize online stores to easily sell fountain pens to a wider audience.
Feed
The feed of a fountain pen is the component that connects the nib of the pen with its ink reservoir. It not only allows the ink to flow to the nib (in what is often described as a "controlled leak") but also regulates the amount of air flowing backwards up to the reservoir to replace this lost ink.
The feed uses a series of narrow channels or "fissures" that run down its lower edge. As ink flows down these fissures, air is simultaneously allowed to flow upwards into the reservoir in an even exchange of volumes. The feed allows ink to flow when the pen is being put to paper but ensures ink does not flow when the pen is not in use. The feed makes use of capillary action; this is noticeable when a pen is refilled with a brightly coloured ink. The ink is taken up and into the feed by way of capillary action (and is often visible in clear demonstrator pens), but is not dispensed onto the paper until the nib makes contact.
How the feed is shaped may determine the wetness and flow of a particular pen. For this reason, feed material alone and its surface roughness may have a significant effect on the way two pens of the same nib size write.
Pen feeds are crucial to preventing ink from dripping or leaking. Feeds often feature finned structures intended for buffering fountain pen ink. Buffering is the capacity to catch and temporarily hold an overflow of ink, caused by conditions other than writing. When a fountain pen nib receives such an overflow it will result in ink blobbing or dripping also known as burping. A pen with a misconfigured feed might fail to deposit any ink whatsoever.
Fiber feeds
Some fountain pens use a fiber wick underneath the nib. They often have a plastic part that looks like a feed that is only used to hold the fiber wick in place and does not assist with ink flow. The mechanism of action is like a felt pen, just with a fountain pen nib on top of it. The fiber feeds offer plenty of ink flow and can stay wet for extended periods. Cleaning fiber feed pens can require longer soaking in water.
Nibs
The modern fountain pen nib is a direct descendant of the iridium-tipped gold dip pen nibs of the 19th century. The earliest attempts at adding a hard and long-wearing tipping material to a gold nib utilized materials such as ruby. A more successful approach exploited the discovery of the platinum group of metals, including ruthenium, osmium, and iridium. From the mid-1830s gold dip pen nibs tipped with iridium were produced in rapidly increasing quantities, first in England and soon thereafter in the United States. The first mass-produced fountain pens used gold nibs sourced from established makers of gold dip pen nibs, some of the most prominent being Mabie Todd, Fairchild, and Aikin Lambert.
Today, nibs are usually made of stainless steel or gold, with the most popular gold alloys being 14 carat (58⅓%) and 18 carat (75%). Titanium is a less common metal used for making nibs. Gold is considered the optimum metal for its flexibility and its resistance to corrosion, although gold's corrosion resistance is less of an issue than in the past because of better stainless steel alloys and less corrosive inks.
Palladium alloys have been used occasionally in the past, usually as a money-saving alternative to white gold. As long as palladium remains more valuable than gold, however, it is unlikely to see much use for nib manufacture.
Nib plating
Further gold plating provides favorable wettability, which is the ability of a solid surface to reduce the surface tension of a liquid in contact with it such that it spreads over the surface.
Nib tipping
Gold and most steel and titanium nibs are tipped with a hard, wear-resistant alloy that typically includes metals from the platinum group. These metals share qualities of extreme hardness and corrosion resistance. The tipping material is often called "iridium", but few if any nib manufacturers have used tipping alloys containing iridium since the mid-1950s. The metals osmium, rhenium, ruthenium, and tungsten are used instead, generally as an alloy, produced as tiny pellets which are welded onto the nib's tip prior to cutting the nib slit and grinding the tip into its final shape. Untipped steel and titanium points will wear more rapidly due to abrasion by the paper.
Capillary action
The nib usually has a tapering or parallel slit cut down its centre, to convey the ink down the nib by capillary action, as well as a "breather hole" of varying shape. The breather hole's intended function is to allow air exchange with the ink reservoir through the channels of the feed, though some modern authorities believe this is a misconception and such venting is unnecessary. Some fountain pens come without a breather hole such as the Camlin Trinity, Monami Olika, Pelikan Pelikano, and Platinum Preppy. The breather hole's other main function is to provide an endpoint to the nib slit and an indexing point for slit cutting. The breather hole also acts as a stress relieving point, preventing the nib from cracking longitudinally from the end of the slit as a result of repeated flexing during use.
The nib narrows to a point where the ink is transferred to the paper. Extremely broad calligraphy pens may have several slits in the nib to increase ink flow and help distribute it evenly across the point, but such designs are more commonly found on dip pens. Nibs divided into three 'tines' are commonly known as music nibs. This is because their line, which can be varied from broad to fine, is suited for writing musical scores.
Types of nibs
Although the most common nibs end in a round point of various sizes (extra fine, fine, medium, broad), various other nib shapes are available. Examples of this are double broad, music, oblique, reverse oblique, stub, italic, and 360-degree nibs.
Broader nibs are used for less precise emphasis, with the benefit of a greater level of ink shading, a property wherein ink pools in parts of a stroke to cause variations in color or sheen – where dyes in ink crystallize on a page instead of absorbing into the paper, which leads to a different color being seen on less absorbent paper due to thin film interference. Finer nibs (e.g. extra fine and fine) may be used for intricate corrections and alterations, at the expense of shading and sheen. Oblique, reverse oblique, stub, and italic nibs may be used for calligraphic purposes or for general handwritten compositions. The line width of a particular nib may vary based on its country of origin; Japanese nibs are often thinner in general.
Nib flexibility
Flexibility is a function of several factors. One is the nib material's resilience; another is its thickness. Finally there is the nib's shape, with longer tines offering more flexibility than short tines, while greater curvature increases stiffness. Contrary to common belief, material alone does not determine a nib's flexibility. Gold alloys of greater purity (18K, or 750/1000 gold) will on average be softer and less springy than alloys of lower purity (14K, or 585/1000 gold), but whatever the alloy its resilience can be altered considerably in manufacture by means of controlled work-hardening.
Fountain pens dating from the first half of the 20th century are more likely to have flexible nibs, suited to the favored handwriting styles of the period (e.g. Copperplate script and Spencerian script). By the 1940s, writing preferences had shifted towards stiffer nibs that could withstand the greater pressure required for writing through carbon paper to create duplicate documents.
Furthermore, competition between the major pen brands such as Parker and Waterman, and the introduction of lifetime guarantees, meant that flexible nibs could no longer be supported profitably. In countries where this rivalry was not present to the same degree, such as the UK and Germany, flexible nibs are more common.
Nowadays, stiff nibs are the norm as people exchange between fountain pens and other writing modes. These more closely emulate the ballpoint pens most modern writers are experienced with. Despite being rigid and firm, the idea that steel nibs write "horribly" is a misconception. More flexible nibs can be easily damaged if excessive pressure is applied to them. Ideally, a fountain pen's nib glides across the paper using the ink as a lubricant, and writing requires no pressure.
Good quality nibs that have been used appropriately are long lasting, often lasting longer than the lifetime of the original owner. Many vintage pens with decades-old nibs can still be used today.
Different nib styles
Other styles of fountain pen nibs include hooded (e.g. Parker 51, Parker 61, 2007 Parker 100, Lamy 2000, and Hero 329), inlaid (e.g. Sheaffer Targa or Sheaffer P.F.M) or integral Nib (Parker T-1, Falcon, and Pilot Myu 701), .
Users are often cautioned not to lend or borrow fountain pens as the nib "wears in" at an angle unique to each individual person. A different user is likely to find that a worn-in nib does not write satisfactorily in their hand and, furthermore, creates a second wear surface, ruining the nib for the original user. This, however, is not a point of concern in pens with modern, durable tipping material, as these pens take many years to develop any significant wear.
Filling mechanisms
Eyedropper filler
The reservoirs of the earliest fountain pens were mostly filled by eyedropper ("dropper" was the term used at the time). This could be messy, spurring development of so-called "self-filling" pens equipped with internal filling mechanisms. Though self-fillers had largely displaced dropper-fillers by the 1920s, they never went out of production, and there has been a revival of interest in recent years. For some, the simplicity, reliability, and large ink capacity of the dropper-filler provide ample compensation for its inconveniences.
After the eyedropper-filler era came the first generation of mass-produced self-fillers, almost all using a rubber sac to hold the ink. The sac was compressed and then released by various mechanisms to fill the pen.
Self-filling designs
The Conklin crescent filler, introduced c. 1901, was one of the first mass-produced self-filling pen designs. The crescent-filling system employs an arch-shaped crescent attached to a rigid metal pressure bar, with the crescent portion protruding from the pen through a slot and the pressure bar inside the barrel. A second component, a C-shaped hard rubber lock ring, is located between the crescent and the barrel. In normal use the ring blocks the crescent from being depressed. To fill the pen, the ring is turned until the gap in the ring is aligned with the crescent, allowing the crescent to be depressed, thus compressing the sac.
Many other variations on the basic sac and pressure bar mechanism were introduced in the first decades of the 20th century, such as the coin-filler (with a slot on the barrel, allowing the pressure bar to be depressed by use of a coin), the matchstick-filler (with a round hole, for a matchstick) and the blow-filler (with a hole at the end of the barrel, through which one blew to compress the sac).
Piston filling innovation
In 1908 Walter A. Sheaffer received a patent for an improved lever-filling pen. Introduced in 1912, Sheaffer's pens sold in rapidly increasing numbers and by 1920 Sheaffer had become one of the largest fountain pen makers in the United States.
Parker introduced the button filler, which had a button hidden beneath a blind cap on the end of the barrel; when pressed, it acted on a pressure bar inside to depress the ink sac.
One of the most complex filling mechanisms was introduced in 1952 with the Sheaffer Snorkel. The Snorkel had an axial tube below the nib that could be extended, allowing the pen to be filled from a bottle without needing to immerse the nib or to wipe it off after filling. With the advent of the modern plastic ink cartridge in the early 1950s most of these filling systems were phased out.
Screw-mechanism piston-fillers were made as early as the 1820s, but the mechanism's modern popularity begins with the original Pelikan of 1929, based upon a patent that was initially licensed to a Croatian company Moster-Penkala by inventor Theodore Kovacs. The basic idea is simple and intuitive: turn a knob at the end of the pen and a screw mechanism draws a piston up the barrel, sucking in ink. Pens with this mechanism remain very popular today. Some of the earlier models had to dedicate as much as half of the pen length to the mechanism.
The advent of telescoping pistons has improved this; the Touchdown Filler was introduced by Sheaffer in 1949. It was advertised as an "Exclusive Pneumatic Down-stroke Filler."
To fill it, a knob at the end of the barrel is unscrewed and the attached plunger is drawn out to its full length. The nib is immersed in ink, the plunger is pushed in, compressing and then releasing the ink sac by means of air pressure. The nib is kept in the ink for approximately 10 seconds to allow the reservoir to fill. This mechanism is very closely modeled after a similar pneumatic filler introduced by Chilton over a decade earlier.
Modern filling mechanisms
A capillary filling system was introduced by Parker in the Parker 61 in 1956. There were no moving parts: the ink reservoir within the barrel was open at the upper end, but contained a tightly rolled length of slotted, flexible plastic. To fill, the barrel was unscrewed, the exposed open end of the reservoir was placed in ink and the interstices of the plastic sheet and slots initiated capillary action, drawing up and retaining the ink. The outside of the reservoir was coated with Teflon, a repellent compound that released excess ink as it was withdrawn. Ink was transferred through a further capillary tube to the nib. No method of flushing the device was offered, and because of problems from clogging with dried and hardened ink, production was eventually stopped.
Around the year 2000, Pelikan introduced a filling system involving a valve in the blind end of the pen, which mates with a specially designed ink bottle. Thus docked, ink is then squeezed into the pen barrel (which, lacking any mechanism other than the valve itself, has nearly the capacity of an eyedropper-fill pen of the same size). This system had been implemented only in their "Level" line, which was discontinued in 2006.
Most pens today use either a piston filler, squeeze-bar filler or cartridge. Many pens are also compatible with a converter, which has the same fitting as the pen's cartridge and has a filling mechanism and a reservoir attached to it. This enables a pen to fill either from cartridges or from a bottle of ink. The most common type of converters are piston-style, but many other varieties may be found today. Piston-style converters generally have a transparent round tubular ink reservoir. Fountain pen inks feature differing surface tensions that can cause an ink to adhere or "stick" against the inside of the reservoir. Common solutions for this problem are adding a small (rust-proof) ink agitating object like a 316L or 904L stainless steel or zirconium dioxide bearing ball, spring or hollow tube in the tubular reservoir to mechanically promote free movement of the contained ink and ink/air exchange during writing. Adding a very small amount of surfactant such as Triton X-100 used in Kodak Photo-Flo 200 wetting agent to the ink will chemically promote free movement of the contained ink and ink/air exchange during writing. However, ink might react adversely to adding a surfactant.
Vacuum fillers, such as those used by Pilot in the Custom 823, utilize air pressure to fill the ink chamber. In this case, while the nib is submerged in ink, a plunger is pushed down the empty chamber to create a vacuum in the space behind it. The end of the chamber has a section wider than the rest, and when the plunger passes this point, the difference in air pressure in the area behind the plunger and the area ahead of it is suddenly evened out and ink rushes in behind the plunger to fill the chamber.
Converters are also available in several different types such as piston, plunger, squeeze and push button in rare cases.
Cartridges
The first commercially successful ink cartridge system for fountain pens was patented in 1890 by the Eagle Pencil Company, using glass cartridges. In the 1920s the John Hancock pen featured cartridges made from thin copper tubing. From the 1930s on, Waterman sold pens in France that used glass cartridges. Cartridge-filling pens only became truly popular in the 1950s, however, with the advent of plastic cartridges. The first was the Italian LUS Atomica in 1952, but it was the Waterman C/F in 1953 that brought cartridge filling to the attention of the international market.
Modern plastic cartridges can contain small ridges on the inside to promote free movement of the contained ink and ink/air exchange during writing. Often cartridges are closed with a small ball that gets pressed into the cartridge during insertion into the pen. This ball also aids free movement of the contained ink.
Standard international
Most European fountain pen brands (for example Conway Stewart, Caran d'Ache, Faber-Castell, Michel Perchin, S.T. Dupont, Montegrappa, Stipula, Pelikan, Montblanc, Europen, Sigma, Delta, Italix, and Rotring) and some pen brands of other continents (for example Acura, Bexley, Esterbrook, Retro51, Tombow, and Platinum (with adaptor)) use so called "international cartridges" (AKA "European cartridges" or "standard cartridges" or "universal cartridges"), in short (38 mm in length, about 0.75 ml of capacity) or long (72 mm, 1.50 ml) sizes, or both. It is to some extent a standard, so the international cartridges of any manufacturer can be used in most fountain pens that accept international cartridges.
Also, converters that are meant to replace international cartridges can be used in most fountain pens that accept international cartridges. Some very compact fountain pens (for example Waterman Ici et La and Monteverde Diva) accept only short international cartridges. Converters can not be used in them (except for so-called mini-converters by Monteverde). Some pens (such as the modern Waterman models) have intentional fittings which prevent the usage of short cartridges. Such pens can only take a proprietary cartridge from the same manufacturer, in this case the long Waterman cartridges.
Proprietary offerings
Many fountain pen manufacturers have developed their own proprietary cartridges, for example Parker, Lamy, Sheaffer, Cross, Sailor, Platinum, Platignum, Waterman, and Namiki. Fountain pens from Aurora, Hero, Duke, and Uranus accept the same cartridges and converters that Parker uses and vice versa (Lamy cartridges, though not officially, are known to interchange with Parker cartridges also). Cartridges of Aurora are slightly different from cartridges by Parker.
Corresponding converters to be used instead of such proprietary cartridges are usually made by the same company that made the fountain pen itself. Some very compact fountain pens accept only proprietary cartridges made by the same company that made that pen, such as Sheaffer Agio Compact and Sheaffer Prelude Compact. It is not possible to use a converter in them at all. In such pens the only practical way to use another brand of ink is to fill empty cartridges with bottled ink using a syringe.
Standard international cartridges are closed by a small ball, held inside the ink exit hole by glue or by a very thin layer of plastic. When the cartridge is pressed into the pen, a small pin pushes in the ball, which falls inside the cartridge. The Parker and Lamy cartridges do not have such a ball. They are closed by a piece of plastic, which is broken by a sharp pin when inserted in the pen.
Concerns and alternatives
Pen manufacturers that produce proprietary cartridges (which in almost all cases are the more expensive ones like the ones mentioned above) tend to discourage the use of cheaper internationally standardised short/long cartridges or adaptations thereof due to their variance in ink quality in the cartridges which may not offer as much performance, or be of lesser quality than the pen's ink cartridge that has been designed specifically for the pen itself. Due to potential additive deficits such as glycerides or uneven pigmentation, cheaper ink may also skip or produce uneven colour on the page and "feather" (spread out of the originally written line) more on thinner grades of paper (e.g. 75 gsm).
While cartridges are more convenient to refill than bottle filling, converter and bottle filling systems are still sold. Non-cartridge filling systems tend to be slightly more economical in the long run since ink is generally less expensive in bottles than in cartridges. Advocates of bottle-based filling systems also cite less waste of plastic for the environment, a wider selection of inks, easier cleaning of pens (as drawing the ink in through the nib helps dissolve old ink), and the ability to check and refill inks at any time.
Inks
Inks intended for use with fountain pens are water-based. These inks are commonly available in bottles. Plastic cartridges came into use in the 1950s, but bottled inks are still the mainstay for many fountain pen enthusiasts. Bottled inks usually cost less than an equivalent amount in cartridges and afford a wider variety of colours and properties.
Fountain pens are not as tightly coupled with their inks as ballpoints or gel pens are, yet some care must be taken when selecting their inks. Contemporary fountain pen inks are almost exclusively dye-based because pigment particles usually clog the narrow passages.
Traditional iron gall inks intended for dip pens are not suitable for fountain pens as they will corrode the pen (a phenomenon known as flash corrosion) and destroy the functionality of the fountain pen. Instead, modern surrogate iron gall formulas are offered for fountain pens. These modern iron gall inks contain a small amount of ferro gallic compounds that are gentler for the inside of a fountain pen, but can still be corrosive if left in the pen for a long period. To avoid corrosion on delicate metal parts and ink clogging, a more thorough than usual cleaning regime – which requires the ink to be flushed out regularly with water – is sometimes advised by manufacturers or resellers.
Some pigmented inks do exist for fountain pens, such as "Carbon Black" and "Chou Kuro" made by the brand Platinum as well as Rohrer & Klingner and De Atramentis inks from Germany, but these are uncommon. Normal India ink cannot be used in fountain pens because it contains shellac as a binder which would very quickly clog such pens.
Inks ideally should be fairly free-flowing, free of sediment, and non-corrosive, though this generally excludes permanence and prevents large-scale commercial use of some colored dyes. Proper care and selection of ink will prevent most problems.
Today
While no longer the primary writing instrument in modern times, fountain pens are used for important official works such as signing valuable documents. Today, fountain pens are often treated as luxury goods and sometimes as status symbols. Fountain pens may serve as an everyday writing instrument, much like the common ballpoint pen. Good quality steel and gold pens are available inexpensively today, particularly in Europe and China, and there are "disposable" fountain pens such as the Pilot Varsity. In France and Germany, in particular, the use of fountain pens is widespread, since in many European countries children still begin to learn writing with this kind of writing device. To avoid mistakes, special ink can be used that can be made invisible by applying an ink eraser.
Fountain pens can serve various artistic purposes such as expressive penmanship and calligraphy, pen and ink artwork, and professional art and design. Many users also favor the perceived elegance, personalization, and sentimentality associated with fountain pens, which computers and ballpoint pens seem to lack, and often state that once they start using fountain pens, ballpoints become awkward to use due to the extra motor effort needed and lack of expressiveness. For ergonomics, fountain pens may relieve physiological stress from writing; alternatives such as the ballpoint pen can induce more pain and damage to those with arthritis. Some also believe they could improve academic performance.
In some countries, fountain pens are usual in lower school grades, believed to teach children better control over writing as many common mistakes of people not used to handwriting (like too much pressure or incorrect hold) feel unnatural or are almost impossible when using traditional pen tips.
Fountain pens are a popular collectable. Fancier pens have been made of precious metals and sometimes even inlaid with jewels. Some are decorated with lacquer, including Japanese maki-e. Avid communities of pen enthusiasts collect and use antique and modern pens and also collect and exchange information about old and modern inks, ink bottles, and inkwells. Collectors may decide to use the antiques in addition to showcasing them in closed spaces such as glass displays.
News outlets report that rather than declining, fountain pen sales have been steadily rising over the last decade. There is a clear resurgence in the appeal and culture of the fountain pen, whether for purposes of collection, enjoyment or as a "lifestyle item". Many agree that the "personal touch" of a fountain pen has led to such a resurgence with modern consumers looking for an alternative in a world of digital products and services.
Amazon reported "sales so far this year [2012] have doubled compared with the same period in 2011. They are four times higher than 2010." The popularity of fountain pens continues to show growth. The market-research firm Euromonitor reported that fountain pen retail sales were up 2.1% in 2016 from a year earlier, reaching $1.046 billion.
| Technology | Writing tools | null |
60288 | https://en.wikipedia.org/wiki/Pewter | Pewter | Pewter () is a malleable metal alloy consisting of tin (85–99%), antimony (approximately 5–10%), copper (2%), bismuth, and sometimes silver. In the past, it was an alloy of tin and lead, but most modern pewter, in order to prevent lead poisoning, is not made with lead. Pewter has a low melting point, around , depending on the exact mixture of metals. The word pewter is possibly a variation of "spelter", a term for zinc alloys (originally a colloquial name for zinc).
History
Pewter was first used around the beginning of the Bronze Age in the Near East. The earliest known piece of pewter was found in an Egyptian tomb, , but it is unlikely that this was the first use of the material. Pewter was used for decorative metal items and tableware in ancient times by the Egyptians and later the Romans, and came into extensive use in Europe from the Middle Ages until the various developments in pottery and glass-making during the 18th and 19th centuries. Pewter was a leading material for producing plates, cups, and bowls before the wide adoption of porcelain. Mass production of pottery, porcelain and glass products have almost universally replaced pewter in daily life, although pewter artifacts continue to be produced, mainly as decorative or specialty items. Pewter was also used around East Asia. Although some items still exist, ancient Roman pewter is rare.
Lidless mugs and lidded tankards may be the most familiar pewter artifacts from the late 17th and 18th centuries, although the metal was also used for many other items including porringers (shallow bowls), plates, dishes, basins, spoons, measures, flagons, communion cups, teapots, sugar bowls, beer steins (tankards), and cream jugs. In the early 19th century, changes in fashion caused a decline in the use of pewter flatware. At the same time, production increased of both cast and spun pewter tea sets, whale-oil lamps, candlesticks, and so on. Later in the century, pewter alloys were often used as a base metal for silver-plated objects.
In the late 19th century, pewter came back into fashion with the revival of medieval objects for decoration. New replicas of medieval pewter objects were created, and collected for decoration. Today, pewter is used in decorative objects, mainly collectible statuettes and figurines, game figures, aircraft and other models, (replica) coins, pendants, plated jewellery and so on. Certain athletic contests, such as the United States Figure Skating Championships, award pewter medals to fourth-place finishers.
Types
In antiquity, pewter was tin alloyed with lead and sometimes also copper. Older pewters with higher lead content are heavier, tarnish faster, and their oxidation has a darker, silver-gray color. Pewters containing lead are no longer used in items that will come in contact with the human body (such as cups, plates, or jewelry), due to the toxicity of lead. Modern pewters are available that are completely free of lead, although many pewters containing lead are still being produced for other purposes.
A typical European casting alloy contains 94% tin, 1% copper and 5% antimony. A European pewter sheet would contain 92% tin, 2% copper, and 6% antimony. Asian pewter, produced mostly in Malaysia, Singapore, and Thailand, contains a higher percentage of tin, usually 97.5% tin, 1% copper, and 1.5% antimony. This makes the alloy slightly softer.
The term Mexican pewter is used for any of various alloys of aluminium that are used for decorative items.
Pewter is also used to imitate platinum in costume jewelry.
Properties
Pewter, being a softer material, can be manipulated in various ways such as being cast, hammered, turned, spun and engraved.
Given that pewter is soft at room temperature, a pewter bell does not ring clearly. Cooling it in liquid nitrogen hardens it and enables it to ring, but also makes it more brittle.
| Physical sciences | Specific alloys | Chemistry |
60343 | https://en.wikipedia.org/wiki/Sediment | Sediment | Sediment is a solid material that is transported to a new location where it is deposited. It occurs naturally and, through the processes of weathering and erosion, is broken down and subsequently transported by the action of wind, water, or ice or by the force of gravity acting on the particles. For example, sand and silt can be carried in suspension in river water and on reaching the sea bed deposited by sedimentation; if buried, they may eventually become sandstone and siltstone (sedimentary rocks) through lithification.
Sediments are most often transported by water (fluvial processes), but also wind (aeolian processes) and glaciers. Beach sands and river channel deposits are examples of fluvial transport and deposition, though sediment also often settles out of slow-moving or standing water in lakes and oceans. Desert sand dunes and loess are examples of aeolian transport and deposition. Glacial moraine deposits and till are ice-transported sediments.
Classification
Sediment can be classified based on its grain size, grain shape, and composition.
Grain size
Sediment size is measured on a log base 2 scale, called the "Phi" scale, which classifies particles by size from "colloid" to "boulder".
Shape
The shape of particles can be defined in terms of three parameters. The form is the overall shape of the particle, with common descriptions being spherical, platy, or rodlike. The roundness is a measure of how sharp grain corners are. This varies from well-rounded grains with smooth corners and edges to poorly rounded grains with sharp corners and edges. Finally, surface texture describes small-scale features such as scratches, pits, or ridges on the surface of the grain.
Form
Form (also called sphericity) is determined by measuring the size of the particle on its major axes. William C. Krumbein proposed formulas for converting these numbers to a single measure of form, such as
where , , and are the long, intermediate, and short axis lengths of the particle. The form varies from 1 for a perfectly spherical particle to very small values for a platelike or rodlike particle.
An alternate measure was proposed by Sneed and Folk:
which, again, varies from 0 to 1 with increasing sphericity.
Roundness
Roundness describes how sharp the edges and corners of particle are. Complex mathematical formulas have been devised for its precise measurement, but these are difficult to apply, and most geologists estimate roundness from comparison charts. Common descriptive terms range from very angular to angular to subangular to subrounded to rounded to very rounded, with increasing degree of roundness.
Surface texture
Surface texture describes the small-scale features of a grain, such as pits, fractures, ridges, and scratches. These are most commonly evaluated on quartz grains, because these retain their surface markings for long periods of time. Surface texture varies from polished to frosted, and can reveal the history of transport of the grain; for example, frosted grains are particularly characteristic of aeolian sediments, transported by wind. Evaluation of these features often requires the use of a scanning electron microscope.
Composition
Composition of sediment can be measured in terms of:
Parent rock lithology
Mineral composition
Chemical make-up
This leads to an ambiguity in which clay can be used as both a size-range and a composition (see clay minerals).
Sediment transport
Sediment is transported based on the strength of the flow that carries it and its own size, volume, density, and shape. Stronger flows will increase the lift and drag on the particle, causing it to rise, while larger or denser particles will be more likely to fall through the flow.
Fluvial processes
Aeolian processes: wind
Wind results in the transportation of fine sediment and the formation of sand dune fields and soils from airborne dust.
Glacial processes
Glaciers carry a wide range of sediment sizes, and deposit it in moraines.
Mass balance
The overall balance between sediment in transport and sediment being deposited on the bed is given by the Exner equation. This expression states that the rate of increase in bed elevation due to deposition is proportional to the amount of sediment that falls out of the flow. This equation is important in that changes in the power of the flow change the ability of the flow to carry sediment, and this is reflected in the patterns of erosion and deposition observed throughout a stream. This can be localized, and simply due to small obstacles; examples are scour holes behind boulders, where flow accelerates, and deposition on the inside of meander bends. Erosion and deposition can also be regional; erosion can occur due to dam removal and base level fall. Deposition can occur due to dam emplacement that causes the river to pool and deposit its entire load, or due to base level rise.
Shores and shallow seas
Seas, oceans, and lakes accumulate sediment over time. The sediment can consist of terrigenous material, which originates on land, but may be deposited in either terrestrial, marine, or lacustrine (lake) environments, or of sediments (often biological) originating in the body of water. Terrigenous material is often supplied by nearby rivers and streams or reworked marine sediment (e.g. sand). In the mid-ocean, the exoskeletons of dead organisms are primarily responsible for sediment accumulation.
Deposited sediments are the source of sedimentary rocks, which can contain fossils of the inhabitants of the body of water that were, upon death, covered by accumulating sediment. Lake bed sediments that have not solidified into rock can be used to determine past climatic conditions.
Key marine depositional environments
The major areas for deposition of sediments in the marine environment include:
Littoral sands (e.g. beach sands, runoff river sands, coastal bars and spits, largely clastic with little faunal content)
The continental shelf (silty clays, increasing marine faunal content).
The shelf margin (low terrigenous supply, mostly calcareous faunal skeletons)
The shelf slope (much more fine-grained silts and clays)
Beds of estuaries with the resultant deposits called "bay mud".
One other depositional environment which is a mixture of fluvial and marine is the turbidite system, which is a major source of sediment to the deep sedimentary and abyssal basins as well as the deep oceanic trenches.
Any depression in a marine environment where sediments accumulate over time is known as a sediment trap.
The null point theory explains how sediment deposition undergoes a hydrodynamic sorting process within the marine environment leading to a seaward fining of sediment grain size.
Environmental issues
Erosion and agricultural sediment delivery to rivers
One cause of high sediment loads is slash and burn and shifting cultivation of tropical forests. When the ground surface is stripped of vegetation and then seared of all living organisms, the upper soils are vulnerable to both wind and water erosion. In a number of regions of the earth, entire sectors of a country have become erodible. For example, on the Madagascar high central plateau, which constitutes approximately ten percent of that country's land area, most of the land area is devegetated, and gullies have eroded into the underlying soil to form distinctive gulleys called lavakas. These are typically wide, long and deep. Some areas have as many as 150 lavakas/square kilometer, and lavakas may account for 84% of all sediments carried off by rivers. This siltation results in discoloration of rivers to a dark red brown color and leads to fish kills. In addition, sedimentation of river basins implies sediment management and siltation costs. The cost of removing an estimated 135 million m3 of accumulated sediments due to water erosion only is likely exceeding 2.3 billion euro (€) annually in the EU and UK, with large regional differences between countries.
Erosion is also an issue in areas of modern farming, where the removal of native vegetation for the cultivation and harvesting of a single type of crop has left the soil unsupported. Many of these regions are near rivers and drainages. Loss of soil due to erosion removes useful farmland, adds to sediment loads, and can help transport anthropogenic fertilizers into the river system, which leads to eutrophication.
The Sediment Delivery Ratio (SDR) is fraction of gross erosion (interill, rill, gully and stream erosion) that is expected to be delivered to the outlet of the river. The sediment transfer and deposition can be modelled with sediment distribution models such as WaTEM/SEDEM. In Europe, according to WaTEM/SEDEM model estimates the Sediment Delivery Ratio is about 15%.
Coastal development and sedimentation near coral reefs
Watershed development near coral reefs is a primary cause of sediment-related coral stress. The stripping of natural vegetation in the watershed for development exposes soil to increased wind and rainfall and, as a result, can cause exposed sediment to become more susceptible to erosion and delivery to the marine environment during rainfall events. Sediment can negatively affect corals in many ways, such as by physically smothering them, abrading their surfaces, causing corals to expend energy during sediment removal, and causing algal blooms that can ultimately lead to less space on the seafloor where juvenile corals (polyps) can settle.
When sediments are introduced into the coastal regions of the ocean, the proportion of land, marine, and organic-derived sediment that characterizes the seafloor near sources of sediment output is altered. In addition, because the source of sediment (i.e., land, ocean, or organically) is often correlated with how coarse or fine sediment grain sizes that characterize an area are on average, grain size distribution of sediment will shift according to the relative input of land (typically fine), marine (typically coarse), and organically-derived (variable with age) sediment. These alterations in marine sediment characterize the amount of sediment suspended in the water column at any given time and sediment-related coral stress.
Biological considerations
In July 2020, marine biologists reported that aerobic microorganisms (mainly), in "quasi-suspended animation", were found in organically-poor sediments, up to 101.5 million years old, 250 feet below the seafloor in the South Pacific Gyre (SPG) ("the deadest spot in the ocean"), and could be the longest-living life forms ever found.
| Physical sciences | Pedology | null |
60372 | https://en.wikipedia.org/wiki/Calamine%20%28mineral%29 | Calamine (mineral) | Calamine is a historic name for an ore of zinc. The name calamine was derived from lapis calaminaris, a Latin corruption of Greek cadmia (καδμία), the old name for zinc ores in general. The name of the Belgian town of Kelmis, La Calamine in French, which was home to a zinc mine, comes from this. In the 18th and 19th centuries large ore mines could be found near the German village of Breinigerberg.
During the early 19th century it was discovered that what had been thought to be one ore was actually two distinct minerals:
Zinc carbonate ZnCO3 or smithsonite and
Zinc silicate Zn4Si2O7(OH)2·H2O or hemimorphite.
Although chemically and crystallographically quite distinct, the two minerals exhibit similar massive or botryoidal external form and are not readily distinguished without detailed chemical or physical analysis. The first person to separate the minerals was the British chemist and mineralogist James Smithson in 1803. In the mining industry the term calamine has been historically used to refer to both minerals indiscriminately.
In mineralogy calamine is no longer considered a valid term. It has been replaced by smithsonite and hemimorphite in order to distinguish it from the pinkish mixture of zinc oxide (ZnO) and iron(III) oxide (Fe2O3) known as calamine lotion.
Early history
In the 16th century demand for latten (brass) in England came from the needs of wool-carding, for which brass-wire combs were preferred, and battery pieces (brassware formed by hammering sheet brass in a battery mill). The only known method for producing the alloy was by heating copper and calamine together in the cementation process and in 1568 a royal charter was granted to the Society of the Mineral and Battery Works to search for the mineral and produce brass, to reduce dependence on imported metal from Germany. Factories to exploit the process were established at Isleworth and Rotherhithe. By the late 17th century enough was known of metallic zinc to make brass solder directly by combining copper and spelter (zinc ingots). In 1738 a patent was granted to William Champion, a Bristol brass founder, for the large-scale reduction of calamine to produce spelter. There were many Calamine mines in Shipham, not far from William Champion's Brass works.
In 1684 a paper presented to the Royal Society described the medicinal and veterinary properties of the compound when in finely powdered form. Since then no mechanism of action for the powder has been identified, and the only medical effect of the powdered mineral appears to be its ability to absorb moisture secreted from irritated and weeping skin.
| Physical sciences | Minerals | Earth science |
14715409 | https://en.wikipedia.org/wiki/Katana | Katana | A is a Japanese sabre characterized by a curved, single-edged blade with a circular or squared guard and long grip to accommodate two hands. Developed later than the tachi, it was used by samurai in feudal Japan and worn with the edge facing upward. Since the Muromachi period, many old tachi were cut from the root and shortened, and the blade at the root was crushed and converted into a katana. The specific term for katana in Japan is and the term katana (刀) often refers to single-edged swords from around the world.
Etymology and loanwords
The word katana first appears in Japanese in the Nihon Shoki of 720. The term is a compound of kata ("one side, one-sided") + na ("blade"), in contrast to the double-sided tsurugi.
The katana belongs to the nihontō family of swords, and is distinguished by a blade length (nagasa) of more than 2 shaku, approximately .
Katana can also be known as dai or daitō among Western sword enthusiasts, although daitō is a generic name for any Japanese long sword, literally meaning "big sword".
As Japanese does not have separate plural and singular forms, both katanas and katana are considered acceptable forms in English.
Pronounced , the kun'yomi (Japanese reading) of the kanji 刀, originally meaning single edged blade (of any length) in Chinese, the word has been adopted as a loanword by the Portuguese. In Portuguese the designation (spelled catana) means "large knife" or machete.
Description
The katana is generally defined as the standard sized, moderately curved (as opposed to the older tachi featuring more curvature) Japanese sword with a blade length greater than 60.6 cm (23.86 inches) (over 2 shaku). It is characterized by its distinctive appearance: a curved, slender, single-edged blade with a circular or squared guard (tsuba) and long grip to accommodate two hands.
With a few exceptions, katana and tachi can be distinguished from each other, if signed, by the location of the signature (mei) on the tang (nakago). In general, the mei should be carved into the side of the nakago which would face outward when the sword was worn. Since a tachi was worn with the cutting edge down, and the katana was worn with the cutting edge up, the mei would be in opposite locations on the tang.
Western historians have said that katana were among the finest cutting weapons in world military history. However, the main weapons on the battlefield in the Sengoku period in the 16th century were yumi (bow), yari (spear), and tanegashima (gun), and katana and tachi were used only for close combat. During this period, the tactics changed to a group battle by ashigaru (foot soldiers) mobilized in large numbers, so naginata and tachi became obsolete as weapons on the battlefield and were replaced by yari and katana. In the relatively peaceful Edo period, katana increased in importance as a weapon, and at the end of the Edo period, shishi (political activists) fought many battles using katana as their main weapon. Katana and tachi were often used as gifts between daimyo (feudal lord) and samurai, or as offerings to the kami enshrined in Shinto shrines, and symbols of authority and spirituality of samurai.
History
The production of swords in Japan is divided into specific time periods:
Jōkotō (ancient swords, until around 900)
Kotō (old swords from around 900–1596)
Shintō (new swords 1596–1780)
Shinshintō (newer swords 1781–1876)
Gendaitō (modern or contemporary swords 1876–present)
Kotō (Old swords)
Katana originates from sasuga (刺刀), a kind of tantō (short sword or knife) used by lower-ranking samurai who fought on foot in the Kamakura period (1185–1333). Their main weapon was a long naginata and sasuga was a spare weapon. In the Nanboku-chō period (1336–1392) which corresponds to the early Muromachi period (1336–1573), long weapons such as ōdachi were popular, and along with this, sasuga lengthened and finally became katana. Also, there is a theory that koshigatana (腰刀), a kind of tantō which was equipped by high ranking samurai together with tachi, developed to katana through the same historical background as sasuga, and it is possible that both developed to katana. The oldest katana in existence today is called Hishizukuri uchigatana, which was forged in the Nanbokuchō period, and was dedicated to Kasuga Shrine later.
The first use of katana as a word to describe a long sword that was different from a tachi, occurs as early as the Kamakura period. These references to "uchigatana" and "tsubagatana" seem to indicate a different style of sword, possibly a less costly sword for lower-ranking warriors. Starting around the year 1400, long swords signed with the katana-style mei were made. This was in response to samurai wearing their tachi in what is now called "katana style" (cutting edge up). Japanese swords are traditionally worn with the mei facing away from the wearer. When a tachi was worn in the style of a katana, with the cutting edge up, the tachi's signature would be facing the wrong way. The fact that swordsmiths started signing swords with a katana signature shows that some samurai of that time period had started wearing their swords in a different manner.
By the 15th century, Japanese swords, including katana, had already gained international fame by being exported to China and Korea. For example, Korea learned how to make Japanese swords by sending swordsmiths to Japan and inviting Japanese swordsmiths to Korea. According to the record of June 1, 1430 in the Veritable Records of the Joseon Dynasty, a Korean swordsmith who went to Japan and mastered the method of making Japanese swords presented a Japanese sword to the King of Korea and was rewarded for the excellent work which was no different from the swords made by the Japanese.
Traditionally, yumi (bows) were the main weapon of war in Japan, and tachi and naginata were used only for close combat. The Ōnin War in the late 15th century in the Muromachi period expanded into a large-scale domestic war, in which employed farmers called ashigaru were mobilized in large numbers. They fought on foot using katana shorter than tachi. In the Sengoku period (period of warring states) in the late Muromachi period, the war became bigger and ashigaru fought in a close formation using yari (spears) lent to them. Furthermore, in the late 16th century, tanegashima (muskets) were introduced from Portugal, and Japanese swordsmiths mass-produced improved products, with ashigaru fighting with leased guns. On the battlefield in Japan, guns and spears became main weapons in addition to bows. Due to the changes in fighting styles in these wars, the tachi and naginata became obsolete among samurai, and the katana, which was easy to carry, became the mainstream. The dazzling looking tachi gradually became a symbol of the authority of high-ranking samurai.
On the other hand, kenjutsu (swordsmanship) that makes use of the characteristics of katana was invented. The quicker draw of the sword was well suited to combat where victory depended heavily on short response times. (The practice and martial art for drawing the sword quickly and responding to a sudden attack was called battōjutsu, which is still kept alive through the teaching of iaido.) The katana further facilitated this by being worn thrust through a belt-like sash (obi) with the sharpened edge facing up. Ideally, samurai could draw the sword and strike the enemy in a single motion. Previously, the curved tachi had been worn with the edge of the blade facing down and suspended from a belt.
From the 15th century, low-quality swords were mass-produced under the influence of the large-scale war. These swords, along with spears, were lent to recruited farmers called ashigaru and swords were exported. Such mass-produced swords are called kazuuchimono, and swordsmiths of the Bisen school and Mino school produced them by division of labor. The export of katana and tachi reached its peak during this period, from the late 15th century to early 16th century when at least 200,000 swords were shipped to Ming dynasty China in official trade in an attempt to soak up the production of Japanese weapons and make it harder for pirates in the area to arm. In the Ming dynasty of China, Japanese swords and their tactics were studied to repel pirates, and wodao and miaodao were developed based on Japanese swords.
From this period, the tang (nakago) of many old tachi were cut and shortened into katana. This kind of remake is called suriage (磨上げ). For example, many of the tachi that Masamune forged during the Kamakura period were converted into katana, so his only existing works are katana and tantō.
From around the 16th century, many Japanese swords were exported to Thailand, where katana-style swords were made and prized for battle and art work, and some of them are in the collections of the Thai royal family.
From the late Muromachi period (Sengoku period) to the early Edo period, samurai were sometimes equipped with a katana blade pointing downwards like a tachi. This style of sword is called handachi, "half tachi". In handachi, both styles were often mixed, for example, fastening to the obi was katana style, but metalworking of the scabbard was tachi style.
In the Muromachi period, especially the Sengoku period, people such as farmers, townspeople, and monks could have a sword. However, in 1588 Toyotomi Hideyoshi banned farmers from owning weapons and conducted a sword hunt to forcibly remove swords from anyone identifying as a farmer.
The length of the katana blade varied considerably during the course of its history. In the late 14th and early 15th centuries, katana blades tended to have lengths between . During the early 16th century, the average length dropped about , approaching closer to . By the late 16th century, the average length had increased again by about , returning to approximately .
Shintō (New swords)
Swords forged after 1596 in the Keichō period of the Azuchi–Momoyama period are classified as shintō (New swords). Japanese swords from shintō are different from kotō in forging method and steel (tamahagane). This is thought to be because Bizen school, which was the largest swordsmith group of Japanese swords, was destroyed by a great flood in 1590 and the mainstream shifted to Mino school, and because Toyotomi Hideyoshi virtually unified Japan, uniform steel began to be distributed throughout Japan. The kotō swords, especially the Bizen school swords made in the Kamakura period, had a midare-utsuri like a white mist between hamon and shinogi, but in the swords from shintō it has almost disappeared. In addition, the whole body of the blade became whitish and hard. Almost no one was able to reproduce midare-utsurii until Kunihira Kawachi reproduced it in 2014.
As the Sengoku period (period of warring states) ended and the Azuchi-Momoyama period to the Edo period started, katana-forging also developed into a highly intricate and well-respected art form. Lacquered saya (scabbards), ornate engraved fittings, silk handles and elegant tsuba (handguards) were popular among samurai in the Edo period, and eventually (especially when Japan was in peace time), katana became more cosmetic and ceremonial items than practical weapons. The Umetada school led by Umetada Myoju who was considered to be the founder of shinto led the improvement of the artistry of Japanese swords in this period. They were both swordsmiths and metalsmiths, and were famous for carving the blade, making metal accouterments such as tsuba (handguard), remodeling from tachi to katana (suriage), and inscriptions inlaid with gold.
During this period, the Tokugawa shogunate required samurai to wear katana and shorter swords in pairs. These short swords were wakizashi and tantō, and wakizashi were mainly selected. This set of two is called a daishō. Only samurai could wear the daishō: it represented their social power and personal honour. Samurai could wear decorative sword mountings in their daily lives, but the Tokugawa shogunate regulated the formal sword that samurai wore when visiting a castle by regulating it as a daisho made of a black scabbard, a hilt wrapped with white ray skin and black string. Japanese swords made in this period are classified as shintō.
Shinshintō (New swords)
In the late 18th century, swordsmith Suishinshi Masahide criticized that the present katana blades only emphasized decoration and had a problem with their toughness. He insisted that the bold and strong kotō blade from the Kamakura period to the Nanboku-chō period was the ideal Japanese sword, and started a movement to restore the production method and apply it to Katana. Katana made after this is classified as a shinshintō. One of the most popular swordsmiths in Japan today is Minamoto Kiyomaro who was active in this shinshintō period. His popularity is due to his timeless exceptional skill, as he was nicknamed "Masamune in Yotsuya" after his disastrous life. His works were traded at high prices and exhibitions were held at museums all over Japan from 2013 to 2014.
The idea that the blade of a sword in the Kamakura period is the best has been continued until now, and as of the 21st century, 80% of Japanese swords designated as National treasure in Japan were made in the Kamakura period, and 70% of them were tachi.
The arrival of Matthew Perry in 1853 and the subsequent Convention of Kanagawa caused chaos in Japanese society. Conflicts began to occur frequently between the forces of sonnō jōi (尊王攘夷派), who wanted to overthrow the Tokugawa Shogunate and rule by the Emperor, and the forces of sabaku (佐幕派), who wanted the Tokugawa Shogunate to continue. These political activists, called the shishi (志士), fought using a practical katana, called the kinnōtō (勤皇刀) or the bakumatsutō (幕末刀). Their katana were often longer than 90 cm (35.43 in) in blade length, less curved, and had a big and sharp point, which was advantageous for stabbing in indoor battles.
Gendaitō (modern or contemporary swords)
Meiji – World War II
During the Meiji period, the samurai class was gradually disbanded, and the special privileges granted to them were taken away, including the right to carry swords in public. The Haitōrei Edict in 1876 forbade the carrying of swords in public except for certain individuals, such as former samurai lords (daimyō), the military, and the police. Skilled swordsmiths had trouble making a living during this period as Japan modernized its military, and many swordsmiths started making other items, such as farm equipment, tools, and cutlery. The craft of making swords was kept alive through the efforts of some individuals, notably Miyamoto Kanenori (宮本包則, 1830–1926) and Gassan Sadakazu (月山貞一, 1836–1918), who were appointed Imperial Household Artist. The businessman Mitsumura Toshimo (光村利藻, 1877－1955) tried to preserve their skills by ordering swords and sword mountings from the swordsmiths and craftsmen. He was especially enthusiastic about collecting sword mountings, and he collected about 3,000 precious sword mountings from the end of the Edo period to the Meiji period. About 1,200 items from a part of his collection are now in the Nezu Museum.
Military action by Japan in China and Russia during the Meiji period helped revive interest in swords, but it was not until the Shōwa period that swords were produced on a large scale again. Japanese military swords produced between 1875 and 1945 are referred to as guntō (military swords).
During the pre-World War II military buildup, and throughout the war, all Japanese officers were required to wear a sword. Traditionally made swords were produced during this period, but in order to supply such large numbers of swords, blacksmiths with little or no knowledge of traditional Japanese sword manufacture were also recruited. In addition, supplies of the Japanese steel (tamahagane) used for swordmaking were limited, so several other types of steel were also used. Quicker methods of forging were also used, such as the use of power hammers, and quenching the blade in oil, rather than hand forging and water. The non-traditionally made swords from this period are called shōwatō, after the regnal name of the Emperor Hirohito, and in 1937, the Japanese government started requiring the use of special stamps on the tang (nakago) to distinguish these swords from traditionally made swords. During this period of war, older antique swords were remounted for use in military mounts. Presently, in Japan, shōwatō are not considered to be "true" Japanese swords, and they can be confiscated. Outside Japan, however, they are collected as historical artifacts.
Post–World War II
Between 1945 and 1953, sword manufacture and sword-related martial arts were banned in Japan. Many swords were confiscated and destroyed, and swordsmiths were not able to make a living. Since 1953, Japanese swordsmiths have been allowed to work, but with severe restrictions: swordsmiths must be licensed and serve a five-year apprenticeship, and only licensed swordsmiths are allowed to produce Japanese swords (nihonto), only two longswords per month are allowed to be produced by each swordsmith, and all swords must be registered with the Japanese Government.
Outside Japan, some of the modern katanas being produced by western swordsmiths use modern steel alloys, such as L6 and A2. These modern swords replicate the size and shape of the Japanese katana and are used by martial artists for iaidō and even for cutting practice (tameshigiri).
Mass-produced swords including iaitō and shinken in the shape of katana are available from many countries, though China dominates the market. These types of swords are typically mass-produced and made with a wide variety of steels and methods.
According to the Parliamentary Association for the Preservation and Promotion of Japanese Swords, organized by Japanese Diet members, many katana distributed around the world as of the 21st century are fake Japanese swords made in China. The Sankei Shimbun analyzed that this is because the Japanese government allowed swordsmiths to make only 24 Japanese swords per person per year in order to maintain the quality of Japanese swords.
Many swordsmiths after the Edo period have tried to reproduce the sword of the Kamakura period which is considered as the best sword in the history of Japanese swords, but they have failed. Then, in 2014, Kunihira Kawachi succeeded in reproducing it and won the Masamune Prize, the highest honor as a swordsmith. No one could win the Masamune Prize unless he made an extraordinary achievement, and in the section of tachi and katana, no one had won for 18 years before Kawauchi.
Types
Katana are distinguished by their type of blade:
Shinogi-Zukuri is the most common blade shape for Japanese katana that provides both speed and cutting power. It features a distinct yokote: a line or bevel that separates the finish of the main blade and the finish of the tip. Shinogi-zukuri was originally produced after the Heian period.
Shobu-Zukuri is a variation of shinogi-zukuri without a yokote, the distinct angle between the long cutting edge and the point section. Instead, the edge curves smoothly and uninterrupted into the point.
Kissaki-Moroha-Zukuri is a katana blade shape with a distinctive curved and double-edged blade. One edge of the blade is shaped in normal katana fashion while the tip is symmetrical and both edges of the blade are sharp.
In addition to these, there are various other types of blades with different shapes, such as Osoraku-zukuri, Unokubi-zukuri, and Kammuri-otoshi-zukuri.
Forging and construction
Typical features of Japanese swords represented by katana and tachi are a three-dimensional cross-sectional shape of an elongated pentagonal to hexagonal blade called shinogi-zukuri, a style in which the blade and the tang (nakago) are integrated and fixed to the hilt (tsuka) with a pin called mekugi, and a gentle curve. When a shinogi-zukuri sword is viewed from the side, there is a ridge line of the thickest part of the blade called shinogi between the cutting-edge side and the back side. This shinogi contributes to lightening and toughening of the blade and high cutting ability.
Katana are traditionally made from a specialized Japanese steel called tamahagane, which is created from a traditional smelting process that results in several, layered steels with different carbon concentrations. This process helps remove impurities and even out the carbon content of the steel. The age of the steel plays a role in the ability to remove impurities, with older steel having a higher oxygen concentration, being more easily stretched and rid of impurities during hammering, resulting in a stronger blade. The smith begins by folding and welding pieces of the steel several times to work out most of the differences in the steel. The resulting block of steel is then drawn out to form a billet.
At this stage, it is only slightly curved or may have no curve at all. The katana's gentle curvature is attained by a process of differential hardening or differential quenching: the smith coats the blade with several layers of a wet clay slurry, which is a special concoction unique to each sword maker, but generally composed of clay, water and any or none of ash, grinding stone powder, or rust. This process is called tsuchioki. The edge of the blade is coated with a thinner layer than the sides and spine of the sword, heated, and then quenched in water (few sword makers use oil to quench the blade). The slurry causes only the blade's edge to be hardened and also causes the blade to curve due to the difference in densities of the micro-structures in the steel. When steel with a carbon content of 0.7% is heated beyond , it enters the austenite phase. When austenite is cooled very suddenly by quenching in water, the structure changes into martensite, which is a very hard form of steel. When austenite is allowed to cool slowly, its structure changes into a mixture of ferrite and pearlite which is softer than martensite.
This process also creates the line down the sides of the blade called the hamon, which is made distinct by polishing. Each hamon and each smith's style of hamon is distinct. Hamon does not refer to the white area on the side of the blade. The white part is the part that is whitened by a polishing process called hadori to make it easier to see the hamon, and the actual hamon is a fuzzy line within the white part. The actual line of the hamon can be seen by holding the sword in your hand and looking at it while changing the angle of the light shining on the blade.
After the blade is forged, it is then sent to be polished. The polishing takes between one and three weeks. The polisher uses a series of successively finer grains of polishing stones in a process called glazing, until the blade has a mirror finish. However, the blunt edge of the katana is often given a matte finish to emphasize the hamon.
Japanese swords are generally made by a division of labor between six and eight craftsmen. Tosho (Toko, Katanakaji) is in charge of forging blades, togishi is in charge of polishing blades, kinkosi (chokinshi) is in charge of making metal fittings for sword fittings, shiroganeshi is in charge of making habaki (blade collar), sayashi is in charge of making scabbards, nurishi is in charge of applying lacquer to scabbards, tsukamakishi is in charge of making hilt, and tsubashi is in charge of making tsuba (hand guard). Tosho use apprentice swordsmiths as assistants. Prior to the Muromachi period, tosho and kacchushi (armorer) used surplus metal to make tsuba, but from the Muromachi period onwards, specialized craftsmen began to make tsuba. Nowadays, kinkoshi sometimes serves as shiroganeshi and tsubashi.
Appreciation
Historically, katana have been regarded not only as weapons but also as works of art, especially for high-quality ones. For a long time, Japanese people have developed a unique appreciation method in which the blade is regarded as the core of their aesthetic evaluation rather than the sword mountings decorated with luxurious lacquer or metal works.
It is said that there are three objects that are the most noteworthy when appreciating a blade. The first is the overall shape referred to as sugata which is the curvature, length, width, tip, and shape of tang of the sword. The second is a fine pattern on the surface of the blade, which is referred to as hada or jigane. By repeatedly folding and forging the blade, fine patterns such as fingerprints, tree rings and bark are formed on its surface. The third is hamon. Hamon is a fuzzy line in the white pattern of the cutting edge produced by quenching and tempering. The object of appreciation is the shape of hamon and the crystal particles formed at the boundary of hamon. Depending on the size of the particles, they can be divided into two types, a nie and a nioi, which makes them look like stars or mist. The pattern, nie and nioi of the hamon are generally difficult to see, and the viewer usually holds the sword in his hand, changing the angle of the light as it hits the blade. In addition to these three objects, a swordsmith signature and a file pattern engraved on tang, and a carving inscribed on the blade, which is referred to as horimono, are also the objects of appreciation.
The Hon'ami clan, which was an authority of appraisal of Japanese swords, rated Japanese swords from these artistic points of view. In addition, experts of modern Japanese swords judge when and by which swordsmith school the sword was made from these artistic points of view.
Generally, the blade and the sword mounting of Japanese swords are displayed separately in museums, and this tendency is remarkable in Japan. For example, the Nagoya Japanese Sword Museum "Nagoya Touken World", one of Japan's largest sword museums, posts separate videos of the blade and the sword mounting on its official website and YouTube.
Rating of Japanese swords and swordsmiths
In Japan, Japanese swords are rated by authorities of each period, and some of the authority of the rating is still valid today.
In 1719, Tokugawa Yoshimune, the 8th shogun of the Tokugawa shogunate, ordered Hon'ami Kōchū, who was an authority of sword appraisal, to record swords possessed by daimyo all over Japan in books. In the completed "Kyōhō Meibutsu Chō" (享保名物帳) 249 precious swords were described, and additional 25 swords were described later. The list also includes 81 swords that had been destroyed in previous fires. The precious swords described in this book were called "Meibutsu" (名物) and the criteria for selection were artistic elements, origins and legends. The list of "Meibutsu" includes 59 swords made by Masamune, 34 by Awataguchi Yoshimitsu and 22 by Go Yoshihiro, and these three swordsmiths were considered special. Daimyo hid some swords for fear that they would be confiscated by the Tokugawa Shogunate, so even some precious swords were not listed in the book. For example, Daihannya Nagamitsu and Yamatorige, which are now designated as National Treasures, were not listed.
Yamada Asaemon V, who was the official sword cutting ability examiner and executioner of the Tokugawa shogunate, published a book "Kaiho Kenjaku" (懐宝剣尺) in 1797 in which he ranked the cutting ability of swords. The book lists 228 swordsmiths, whose forged swords are called "Wazamono" (業物) and the highest "Saijo Ō Wazamono" (最上大業物) has 12 selected. In the reprinting in 1805, one swordsmith was added to the highest grade, and in the major revised edition in 1830 "Kokon Kajibiko" (古今鍛冶備考), two swordsmiths were added to the highest grade, and in the end, 15 swordsmiths were ranked as the highest grade. The katana forged by Nagasone Kotetsu, one of the top-rated swordsmith, became very popular at the time when the book was published, and many counterfeits were made. In these books, the three swordsmiths treated specially in "Kyōhō Meibutsu Chō" and Muramasa, who was famous at that time for forging swords with high cutting ability, were not mentioned. The reasons for this are considered to be that Yamada was afraid of challenging the authority of the shogun, that he could not use the precious sword possessed by the daimyo in the examination, and that he was considerate of the legend of Muramasa's curse.
At present, by the Law for the Protection of Cultural Properties, important swords of high historical value are designated as Important Cultural Properties (Jūyō Bunkazai, 重要文化財), and special swords among them are designated as National Treasures (Kokuhō, 国宝). The swords designated as cultural properties based on the law of 1930, which was already abolished, have the rank next to Important Cultural Properties as Important Art Object (Jūyō Bijutsuhin, 重要美術品). In addition, the Society for Preservation of Japanese Art Swords, a public interest incorporated foundation, rates high-value swords in four grades, and the highest grade Special Important Sword (Tokubetsu Juyo Token, 特別重要刀剣) is considered to be equivalent to the value of Important Art Object. Although swords owned by the Japanese imperial family are not designated as National Treasures or Important Cultural Properties because they are outside the jurisdiction of the Law for the Protection of Cultural Properties, there are many swords of the National Treasure class, and they are called "Gyobutsu" (御物).
Currently, there are several authoritative rating systems for swordsmiths. According to the rating approved by the Japanese government, from 1890 to 1947, two swordsmiths who were appointed as Imperial Household Artist and after 1955, six swordsmiths who were designated as Living National Treasure are regarded as the best swordsmiths. According to the rating approved by The Society for Preservation of Japanese Art Swords, a public interest incorporated foundation, 39 swordsmiths who were designated as Mukansa (無鑑査) since 1958 are considered to be the highest ranking swordsmiths. The best sword forged by Japanese swordsmiths is awarded the most honorable Masamune prize by The Society for Preservation of Japanese Art Swords. Since 1961, eight swordsmiths have received the Masamune Prize, and among them, three swordsmiths, Masamine Sumitani, Akitsugu Amata and Toshihira Osumi, have received the prize three times each and Sadakazu Gassan II has received the prize two times. These four people were designated both Living National Treasures and Mukansa.
Usage in martial arts
Katana were used by samurai both in the battlefield and for practicing several martial arts, and modern martial artists still use a variety of katana. Martial arts in which training with katana is used include aikidō, iaijutsu, battōjutsu, iaidō, kenjutsu, kendō, ninjutsu,Tenshin Shōden Katori Shintō-ryū and Shinkendo. However, for safety reasons, katana used for martial arts are usually blunt edged iaito or wooden bokken, to reduce the risk of injury. Sharp katana are only really used during tameshigiri (blade testing), where a practitioner practices cutting a bamboo or tatami straw post.
Storage and maintenance
If mishandled in its storage or maintenance, the katana may become irreparably damaged. The blade should be stored horizontally in its sheath, curve down and edge facing upward to maintain the edge. It is extremely important that the blade remain well-oiled, powdered and polished, as the natural moisture residue from the hands of the user will rapidly cause the blade to rust if not cleaned off. The traditional oil used is chōji oil (99% mineral oil and 1% clove oil for fragrance). Similarly, when stored for longer periods, it is important that the katana be inspected frequently and aired out if necessary in order to prevent rust or mold from forming (mold may feed off the salts in the oil used to polish the blade).
World records
Multiple sword world records were made with a katana and verified by Guinness World Records. Iaido master Isao Machii set the record for "Most martial arts katana cuts to one mat (suegiri)", "Fastest 1,000 martial arts sword cuts", "Most sword cuts to straw mats in three minutes", and "Fastest tennis ball (708km/h) cut by sword". There are various records for Tameshigiri. For example, the Greek Agisilaos Vesexidis set the record for most martial arts sword cuts in one minute (73) on 25 June 2016.
Ownership and trade restrictions
Republic of Ireland
Under the Firearms and Offensive Weapons Act 1990 (Offensive Weapons) (Amendment) Order 2009, katanas made post-1953 are illegal unless made by hand according to traditional methods.
United Kingdom
As of April 2008, the British government added swords with a curved blade of 50 cm (20 in) or over in length ("the length of the blade shall be the straight line distance from the top of the handle to the tip of the blade") to the Offensive Weapons Order. This ban was a response to reports that samurai swords were used in more than 80 attacks and four killings over the preceding four years. Those who violate the ban would be jailed up to six months and charged a fine of £5,000. Martial arts practitioners, historical re-enactors and others may still own such swords. The sword can also be legal provided it was made in Japan before 1954, or was made using traditional sword making methods. It is also legal to buy if it can be classed as a "martial artist's weapon". This ban applies to England, Wales, Scotland and Northern Ireland. This ban was amended in August 2008 to allow sale and ownership without licence of "traditional" hand-forged katana.
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| Technology | Melee weapons | null |
3110432 | https://en.wikipedia.org/wiki/Pi%20Puppis | Pi Puppis | Pi Puppis, Latinized from π Puppis, also named Ahadi, is the second-brightest star in the southern constellation of Puppis. It has an apparent visual magnitude of 2.733, so it can be viewed with the naked eye at night. Parallax measurements yield an estimated distance of roughly from the Earth. This is a double star with a magnitude 6.86 companion at an angular separation of 0.72 arcsecond and a position angle of 148° from the brighter primary.
The spectrum of Pi Puppis matches a stellar classification of K3 Ib. The Ib luminosity class indicates this a lower luminosity supergiant star that has consumed the hydrogen fuel at its core, evolved away from the main sequence, and expanded to about 235 times the Sun's radius. The effective temperature of the star's outer envelope is approximately 4,000 K, which gives it the orange hue of a K-type star. With a mass 11.7 times that of the Sun, this is a short-lived star with an estimated age of 20 million years.
It is a semiregular variable star that varies in apparent magnitude from a high of 2.70 down to 2.85. Pi Puppis is the brightest star in the open cluster Collinder 135.
Naming
The star has the traditional name Ahadi, which is derived from Arabic for "having much promise". In Chinese, (), meaning Bow and Arrow, refers to an asterism consisting of π Puppis, δ Canis Majoris, η Canis Majoris, HD 63032, HD 65456, ο Puppis, k Puppis, ε Canis Majoris and κ Canis Majoris. Consequently, π Puppis itself is known as (, .)
| Physical sciences | Notable stars | Astronomy |
3112575 | https://en.wikipedia.org/wiki/Intertidal%20zone | Intertidal zone | The intertidal zone or foreshore is the area above water level at low tide and underwater at high tide; in other words, it is the part of the littoral zone within the tidal range. This area can include several types of habitats with various species of life, such as sea stars, sea urchins, and many species of coral with regional differences in biodiversity. Sometimes it is referred to as the littoral zone or seashore, although those can be defined as a wider region.
The intertidal zone also includes steep rocky cliffs, sandy beaches, bogs or wetlands (e.g., vast mudflats). This area can be a narrow strip, such as in Pacific islands that have only a narrow tidal range, or can include many meters of shoreline where shallow beach slopes interact with high tidal excursion. The peritidal zone is similar but somewhat wider, extending from above the highest tide level to below the lowest. Organisms in the intertidal zone are well-adapted to their environment, facing high levels of interspecific competition and the rapidly changing conditions that come with the tides. The intertidal zone is also home to several species from many different phyla (Porifera, Annelida, Coelenterata, Mollusca, Arthropoda, etc.).
The water that comes with the tides can vary from brackish waters, fresh with rain, to highly saline and dry salt, with drying between tidal inundations. Wave splash can dislodge residents from the littoral zone. With the intertidal zone's high exposure to sunlight, the temperature can range from very hot with full sunshine to near freezing in colder climates. Some microclimates in the littoral zone are moderated by local features and larger plants such as mangroves. Adaptations in the littoral zone allow the utilization of nutrients supplied in high volume on a regular basis from the sea, which is actively moved to the zone by tides. The edges of habitats, in this case the land and sea, are themselves often significant ecosystems, and the littoral zone is a prime example.
A typical rocky shore can be divided into a spray zone or splash zone (also known as the supratidal zone), which is above the spring high-tide line and is covered by water only during storms, and an intertidal zone, which lies between the high and low tidal extremes. Along most shores, the intertidal zone can be clearly separated into the following subzones: high tide zone, middle tide zone, and low tide zone. The intertidal zone is one of a number of marine biomes or habitats, including estuaries, the neritic zone, the photic zone, and deep zones.
Zonation
Marine biologists divide the intertidal region into three zones (low, middle, and high), based on the overall average exposure of the zone. The low intertidal zone, which borders on the shallow subtidal zone, is only exposed to air at the lowest of low tides and is primarily marine in character. The mid intertidal zone is regularly exposed and submerged by average tides. The high intertidal zone is only covered by the highest of the high tides, and spends much of its time as terrestrial habitat. The high intertidal zone borders on the splash zone (the region above the highest still-tide level, but which receives wave splash). On shores exposed to heavy wave action, the intertidal zone will be influenced by waves, as the spray from breaking waves will extend the intertidal zone.
Depending on the substratum and topography of the shore, additional features may be noticed. On rocky shores, tide pools form in depressions that fill with water as the tide rises. Under certain conditions, such as those at Morecambe Bay, quicksand may form.
Low tide zone (lower littoral)
This subregion is mostly submerged – it is only exposed at the point of low tide and for a longer period of time during extremely low tides. This area is teeming with life; the most notable difference between this subregion and the other three is that there is much more marine vegetation, especially seaweeds. There is also a great biodiversity. Organisms in this zone generally are not well adapted to periods of dryness and temperature extremes. Some of the organisms in this area are abalone, sea anemones, brown seaweed, chitons, crabs, green algae, hydroids, isopods, limpets, mussels, nudibranchs, sculpin, sea cucumber, sea lettuce, sea palms, starfish, sea urchins, shrimp, snails, sponges, surf grass, tube worms, and whelks. Creatures in this area can grow to larger sizes because there is more available energy in the localized ecosystem. Also, marine vegetation can grow to much greater sizes than in the other three intertidal subregions due to the better water coverage. The water is shallow enough to allow plenty of sunlight to reach the vegetation to allow substantial photosynthetic activity, and the salinity is at almost normal levels. This area is also protected from large predators such as fish because of the wave action and the relatively shallow water.
Ecology
The intertidal region is an important model system for the study of ecology, especially on wave-swept rocky shores. The region contains a high diversity of species, and the zonation created by the tides causes species ranges to be compressed into very narrow bands. This makes it relatively simple to study species across their entire cross-shore range, something that can be extremely difficult in, for instance, terrestrial habitats that can stretch thousands of kilometres. Communities on wave-swept shores also have high turnover due to disturbance, so it is possible to watch ecological succession over years rather than decades.
The burrowing invertebrates that make up large portions of sandy beach ecosystems are known to travel relatively great distances in cross-shore directions as beaches change on the order of days, semilunar cycles, seasons, or years. The distribution of some species has been found to correlate strongly with geomorphic datums such as the high tide strand and the water table outcrop.
Since the foreshore is alternately covered by the sea and exposed to the air, organisms living in this environment must be adapted to both wet and dry conditions. Intertidal zone biomass reduces the risk of shoreline erosion from high intensity waves. Typical inhabitants of the intertidal rocky shore include sea urchins, sea anemones, barnacles, chitons, crabs, isopods, mussels, starfish, and many marine gastropod molluscs such as limpets and whelks. Sexual and asexual reproduction varies by inhabitants of the intertidal zones.
Humans have historically used intertidal zones as foraged food sources during low tide . Migratory birds also rely on intertidal species for feeding areas because of low water habitats consisting of an abundance of mollusks and other marine species.
Legal issues
As with the dry sand part of a beach, legal and political disputes can arise over the ownership and use of the foreshore. One recent example is the New Zealand foreshore and seabed controversy. In legal discussions, the foreshore is often referred to as the wet-sand area.
For privately owned beaches in the United States, some states such as Massachusetts use the low-water mark as the dividing line between the property of the State and that of the beach owner; however the public still has fishing, fowling, and navigation rights to the zone between low and high water. Other states such as California use the high-water mark.
In the United Kingdom, the foreshore is generally deemed to be owned by the Crown, with exceptions for what are termed several fisheries, which can be historic deeds to title, dating back to King John's time or earlier, and the Udal Law, which applies generally in Orkney and Shetland.
In Greece, according to the L. 2971/01, the foreshore zone is defined as the area of the coast that might be reached by the maximum climbing of the waves on the coast (maximum wave run-up on the coast) in their maximum capacity (maximum referring to the "usually maximum winter waves" and of course not to exceptional cases, such as tsunamis). The foreshore zone, a part of the exceptions of the law, is public, and permanent constructions are not allowed on it. In Italy, about half the shoreline is owned by the government but leased to private beach clubs called lidos.
In the East African and West Indian Ocean, intertidal zone management is often neglected of being a priority due to there being no intent for collective economic productivity. According to workshops performing questionaries, it is stated that eighty-six percent of respondents believe mismanagement of mangrove and coastal ecosystems are due to lack of knowledge to steward the ecosystems, yet forty-four percent of respondents state that there is a fair amount of knowledge used in those regions for fisheries.
Threats
Intertidal zones are sensitive habitats with an abundance of marine species that can experience ecological hazards associated with tourism and human-induced environmental impacts. A variety of other threats that have been summarized by scientists include nutrient pollution, overharvesting, habitat destruction, and climate change. Habitat destruction is advanced through activities including harvesting fisheries with drag nets and a neglect of the sensitivity of intertidal zones.
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| Physical sciences | Oceanography | Earth science |
17557798 | https://en.wikipedia.org/wiki/Heat%20illness | Heat illness | Heat illness is a spectrum of disorders due to increased body temperature. It can be caused by either environmental conditions or by exertion. It includes minor conditions such as heat cramps, heat syncope, and heat exhaustion as well as the more severe condition known as heat stroke. It can affect any or all anatomical systems. Heat illnesses include: heat stroke, heat exhaustion, heat syncope, heat edema, heat cramps, heat rash, heat tetany.
Prevention includes avoiding medications that can increase the risk of heat illness, gradual adjustment to heat, and sufficient fluids and electrolytes.
Classification
A number of heat illnesses exist including:
Heat stroke - Defined by a body temperature of greater than due to environmental heat exposure with lack of thermoregulation. Symptoms include dry skin, rapid, strong pulse and dizziness.
Heat exhaustion - Can be a precursor of heatstroke; the symptoms include heavy sweating, rapid breathing and a fast, weak pulse.
Heat syncope - Fainting or dizziness as a result of overheating.
Heat edema - Swelling of extremities due to water retention following dilation of blood vessels in response to heat.
Heat cramps - Muscle pains that happen during heavy exercise in hot weather.
Heat rash - Skin irritation from excessive sweating.
Heat tetany - Usually results from short periods of stress in intense heat. Symptoms may include hyperventilation, respiratory problems, numbness or tingling, or muscle spasms.
Overview of diseases
Hyperthermia, also known as heat stroke, becomes commonplace during periods of sustained high temperature and humidity. Older adults, very young children, and those who are sick or overweight are at a higher risk for heat-related illness. The chronically ill and elderly are often taking prescription medications (e.g., diuretics, anticholinergics, antipsychotics, and antihypertensives) that interfere with the body's ability to dissipate heat.
Heat edema presents as a transient swelling of the hands, feet, and ankles and is generally secondary to increased aldosterone secretion, which enhances water retention. When combined with peripheral vasodilation and venous stasis, the excess fluid accumulates in the dependent areas of the extremities. The heat edema usually resolves within several days after the patient becomes acclimated to the warmer environment. No treatment is required, although wearing support stockings and elevating the affected legs will help minimize the edema.
Heat rash, also known as prickly heat, is a maculopapular rash accompanied by acute inflammation and blocked sweat ducts. The sweat ducts may become dilated and may eventually rupture, producing small pruritic vesicles on an erythematous base. Heat rash affects areas of the body covered by tight clothing. If this continues for a duration of time it can lead to the development of chronic dermatitis or a secondary bacterial infection. Prevention is the best therapy. It is also advised to wear loose-fitting clothing in the heat. Once heat rash has developed, the initial treatment involves the application of chlorhexidine lotion to remove any desquamated skin. The associated itching may be treated with topical or systemic antihistamines. If infection occurs a regimen of antibiotics is required.
Heat cramps are painful, often severe, involuntary spasms of the large muscle groups used in strenuous exercise. Heat cramps tend to occur after intense exertion. They usually develop in people performing heavy exercise while sweating profusely and replenishing fluid loss with non-electrolyte containing water. This is believed to lead to hyponatremia that induces cramping in stressed muscles. Rehydration with salt-containing fluids provides rapid relief. Patients with mild cramps can be given oral .2% salt solutions, while those with severe cramps require IV isotonic fluids. The many sport drinks on the market are a good source of electrolytes and are readily accessible.
Heat syncope is related to heat exposure that produces orthostatic hypotension. This hypotension can precipitate a near-syncopal episode. Heat syncope is believed to result from intense sweating, which leads to dehydration, followed by peripheral vasodilation and reduced venous blood return in the face of decreased vasomotor control. Management of heat syncope consists of cooling and rehydration of the patient using oral rehydration therapy (sport drinks) or isotonic IV fluids. People who experience heat syncope should avoid standing in the heat for long periods of time. They should move to a cooler environment and lie down if they recognize the initial symptoms. Wearing support stockings and engaging in deep knee-bending movements can help promote venous blood return.
Heat exhaustion is considered by experts to be the forerunner of heat stroke (hyperthermia). It may even resemble heat stroke, with the difference being that the neurologic function remains intact. Heat exhaustion is marked by excessive dehydration and electrolyte depletion. Symptoms may include diarrhea, headache, nausea and vomiting, dizziness, tachycardia, malaise, and myalgia. Definitive therapy includes removing patients from the heat and replenishing their fluids. Most patients will require fluid replacement with IV isotonic fluids at first. The salt content is adjusted as necessary once the electrolyte levels are known. After discharge from the hospital, patients are instructed to rest, drink plenty of fluids for 2–3 hours, and avoid the heat for several days. If this advice is not followed it may then lead to heat stroke.
Symptoms
Increased temperatures have been reported to cause heat stroke, heat exhaustion, heat syncope, and heat cramps. Some studies have also looked at how severe heat stroke can lead to permanent damage to organ systems. This damage can increase the risk of early mortality because the damage can cause severe impairment in organ function. Other complications of heat stroke include respiratory distress syndrome in adults and disseminated intravascular coagulation. Some researchers have noted that any compromise to the human body's ability to thermoregulate would in theory increase risk of mortality. This includes illnesses that may affect a person's mobility, awareness, or behavior.
Prevention
Prevention includes avoiding medications that can increase the risk of heat illness (e.g. antihypertensives, diuretics, and anticholinergics), gradual adjustment to heat, and sufficient fluids and electrolytes.
Some common medications that have an effect on thermoregulation can also increase the risk of mortality. Specific examples include anticholinergics, diuretics, phenothiazines and barbiturates.
Epidemiology
Heat stroke is relatively common in sports. About 2 percent of sports-related deaths that occurred in the United States between 1980 and 2006 were caused by exertional heat stroke. Football in the United States has the highest rates. The month of August, which is associated with pre-season football camps across the country, accounts for 66.3% of exertion heat-related illness time-loss events. Heat illness is also not limited geographically and is widely distributed throughout the United States. An average of 5946 persons were treated annually in US hospital emergency departments (2 visits/ 100,00 population) with a hospitalization rate of 7.1%. Most commonly males are brought in 72.5% and persons 15–19 years of age 35.6% When taking into consideration all high school athletes, heat illness occurs at a rate of 1.2 per 100,000 kids. When comparing risk by sport, Football was 11.4 times more likely than all other sports combined to be exposed to an exertional heat illness.
Between 1999 and 2003, the US had a total of 3442 deaths from heat illness. Those who work outdoors are at particular risk for heat illness, though those who work in poorly-cooled spaces indoors are also at risk. Between 1992 and 2006, 423 workers died from heat illness in the US. Exposure to environmental heat led to 37 work-related deaths. There were 2,830 nonfatal occupational injuries and illnesses involving days away from work as well, in 2015. Kansas had the highest heat related injury while on the job with a rate of 1.3 per 10,000 workers, while Texas had the most overall. Due to the much higher state population of Texas, their prevalence was only 0.4 per 10,000 or 4 per 100,000. Of the 37 deaths reported heat illnesses, 33 of the 37 occurred between the summer months of June through September. The most dangerous profession that was documented was transportation and material moving. Transportation and material moving accounted for 720 of the 2,830 reported nonfatal occupational injuries or 25.4 percent. After transportation and material moving, Production placed second followed by protective services, installation, maintenance, and repair and construction all in succession
Effects of climate change
A 2016 U.S. government report said that climate change could result in "tens of thousands of additional premature deaths per year across the United States by the end of this century." Indeed, between 2014 and 2017, heat exposure deaths tripled in Arizona (76 deaths in 2014; 235 deaths in 2017) and increased fivefold in Nevada (29 deaths in 2014; 139 deaths in 2017).
History
Heat illness used to be blamed on a tropical fever named calenture.
| Biology and health sciences | Types | Health |
617607 | https://en.wikipedia.org/wiki/Psilomelane | Psilomelane | Psilomelane is a group name for hard black manganese oxides including hollandite and romanechite. Psilomelane consists of hydrous manganese oxide with variable amounts of barium and potassium. Psilomelane is erroneously, and uncommonly, known as black hematite, despite not being related to true hematite, which is an iron oxide.
Formula
Generalized formula may be represented as or as . It is sometimes considered to be a hydrous manganese manganate, but of doubtful composition. The amount of manganese present corresponds to 70-80% of manganous oxide with 10-15% of available oxygen.
Characteristics
Psilomelane has no definite chemical composition and occurs as botryoidal and stalactitic masses with a smooth shining surface and submetallic lustre. The mineral is readily distinguished from other hydrous manganese oxides (manganite and wad) by its greater hardness 5 to 6; the specific gravity varies from 3.7 to 4.7. The streak is brownish black and the fracture smooth. The mineral often contains admixed impurities, such as iron hydroxides. It is soluble in hydrochloric acid with evolution of chlorine gas.
History and occurrence
The name, dating back to 1758, makes reference to its characteristic appearance, from the ancient Greek ψιλός: psilos for (naked, smooth, bald) and μέλας: melas (black); a Latinized form is calvonigrite from calvo for (bald, smooth) and negri (black).
It is a common and important ore of manganese, occurring under the same conditions and having the same commercial applications as pyrolusite. It is found at many localities; amongst those which have yielded typical botryoidal specimens may be mentioned the Restormel iron mine at Lostwithiel in Cornwall, Brendon Hills in Somerset, Hoy in Orkney, Sayn near Coblenz, Hout Bay near Cape Town, and Crimora in Augusta county, Virginia. With pyrolusite it is extensively mined in Vermont, Virginia, Arkansas, and Nova Scotia.
Image gallery
| Physical sciences | Minerals | Earth science |
618063 | https://en.wikipedia.org/wiki/Satellite%20phone | Satellite phone | A satellite telephone, satellite phone or satphone is a type of mobile phone that connects to other phones or the telephone network by radio link through satellites orbiting the Earth instead of terrestrial cell sites, as cellphones do. Therefore, they can work in most geographic locations on the Earth's surface, as long as open sky and the line-of-sight between the phone and the satellite are provided. Depending on the architecture of a particular system, coverage may include the entire Earth or only specific regions. Satellite phones provide similar functionality to terrestrial mobile telephones; voice calling, text messaging, and low-bandwidth Internet access are supported through most systems. The advantage of a satellite phone is that it can be used in such regions where local terrestrial communication infrastructures, such as landline and cellular networks, are not available.
Satellite phones are popular on expeditions into remote locations where there is no reliable cellular service, such as recreational hiking, hunting, fishing, and boating trips, as well as for business purposes, such as mining locations and maritime shipping. Satellite phones rarely get disrupted by natural disasters on Earth or human actions such as war, so they have proven to be dependable communication tools in emergency and humanitarian situations, when the local communications system have been compromised.
The mobile equipment, also known as a terminal, varies widely. Early satellite phone handsets had a size and weight comparable to that of a late-1980s or early-1990s mobile phone, but usually with a large retractable antenna. More recent satellite phones are similar in size to a regular mobile phone while some prototype satellite phones have no distinguishable difference from an ordinary smartphone.
A fixed installation, such as one used aboard a ship, may include large, rugged, rack-mounted electronics, and a steerable microwave antenna on the mast that automatically tracks the overhead satellites. Smaller installations using VoIP over a two-way satellite broadband service such as BGAN or VSAT bring the costs within the reach of leisure vessel owners. Internet service satellite phones have notoriously poor reception indoors, though it may be possible to get a consistent signal near a window or in the top floor of a building if the roof is sufficiently thin. The phones have connectors for external antennas that can be installed in vehicles and buildings. The systems also allow for the use of repeaters, much like terrestrial mobile phone systems.
In the early 2020s, various manufacturers began to integrate satellite messaging connectivity and satellite emergency services into conventional mobile phones for use in remote regions, where there is no reliable terrestrial network.
History
The first satellite relayed phone calls were achieved early on in the space age, after the first relay test was conducted by Pioneer 1 and the first broadcast by SCORE in 1958 at the end of the year, after Sputnik I became at the beginning of the year the first satellite in history.
MARISAT was the first mobile communications satellite, eventually operated by the first privatized satellite communication INMARSAT organization, which was formed in 1979.
Satellite network
Satellite phone systems can be classified into two types: systems that use satellites in a high geostationary orbit, above the Earth's surface, and systems that use satellites in low Earth orbit (LEO), above the Earth.
Geostationary satellites
Some satellite phones use satellites in geostationary orbit (GSO), which appear at a fixed position in the sky. These systems can maintain near-continuous global coverage with only three or four satellites, reducing the launch costs. The satellites used for these systems are very heavy (about 5000 kg) and expensive to build and launch. The satellites orbit at an altitude of above the Earth's surface; a noticeable delay is present while making a phone call or using data services due to the large distance from users. The amount of bandwidth available on these systems is substantially higher than that of the low Earth orbit systems; all three active systems provide portable satellite Internet using laptop-sized terminals with speeds ranging from 60 to 512 kbit per second (kbps).
Geostationary satellite phones can only be used at lower latitudes, generally between 70 degrees north of the equator and 70 degrees south of the equator. At higher latitudes the satellite appears at such a low angle in the sky that radio frequency interference from terrestrial sources in the same frequency bands can interfere with the signal.
Another disadvantage of geostationary satellite systems is that in many areas—even where a large amount of open sky is present—the line-of-sight between the phone and the satellite is broken by obstacles such as steep hills and forest. The user will need to find an area with line-of-sight before using the phone. This is not the case with LEO services: even if the signal is blocked by an obstacle, one can wait a few minutes until another satellite passes overhead, but a GSO satellite may drop a call when line of sight is lost.
ACeS: This former Indonesia-based small regional operator provided voice and data services in East Asia, South Asia, and Southeast Asia using a single satellite. It ceased operations in 2014.
Inmarsat: The oldest satellite phone operator, a British company founded in 1979. It originally provided large fixed installations for ships, but has recently entered the market of hand-held phones in a joint venture with ACeS. The company operates eleven satellites. Coverage is available on most of the Earth, except polar regions.
Thuraya: Established in 1997, United Arab Emirates-based Thuraya's satellites provide coverage across Europe, Africa, the Middle East, Asia and Australia.
MSAT / SkyTerra: An American satellite-phone company that uses equipment similar to Inmarsat, but plans to launch a service using hand-held devices in the Americas similar to Thuraya's.
Terrestar: Satellite-phone system for North America.
ICO Global Communications: An American satellite-phone company which has launched a single geosynchronous satellite, not yet active.
Tiantong: A Chinese satellite-phone system is planned to provide satellite phone and short message sending and receiving functions to users in China and its surrounding areas, the Middle East, Africa and other related regions, as well as most of the Pacific and Indian Oceans.
Low Earth orbit
Satellite phones may utilize satellites in low Earth orbit (LEO). The advantages include the possibility of providing worldwide wireless coverage with no gaps. LEO satellites orbit the Earth in high-speed, low-altitude orbits with an orbital time of 70–100 minutes, an altitude of . Since the satellites are not geostationary, they move with respect to the ground. Any given satellite is only in view of a phone for a short time, so the call must be "handed off" electronically to another satellite when one passes beyond the local horizon. Depending on the positions of both the satellite and terminal, a usable pass of an individual LEO satellite will typically last 4–15 minutes on average. At least one satellite must have line-of-sight to every coverage area at all times to guarantee coverage; thus a constellation of satellites, typically 40 to 70, is required to maintain worldwide coverage.
Globalstar: A network covering most of the world's landmass using 48 active satellites. However, many areas of the Earth's surface are left without coverage, since a satellite requires to be in range of an Earth station gateway. Satellites fly in an inclined orbit of 52 degrees, therefore polar regions cannot be covered. The network went into full commercial service in February 2000. A second-generation constellation consists of 24 low Earth orbiting (LEO) satellites. The launch of the second-generation constellation was completed on February 6, 2013.
Iridium: A network operating 66 satellites in a polar orbit that claims to have coverage everywhere on Earth. Radio cross-links are used between satellites to relay data to the nearest satellite with a connection to an Earth station. Commercial service started in November 1998 and fell into Chapter 11 bankruptcy in August 1999. In 2001, service was re-established by Iridium Satellite LLC. Iridium NEXT, a second-generation constellation of the communications satellites, was completed on January 11, 2019.
Both systems, based in the United States, started in the late 1990s, but soon went into bankruptcy after failing to gain enough subscribers to fund launch costs. They are now operated by new owners who bought the assets for a fraction of their original cost and are now both planning to launch replacement constellations supporting higher bandwidth. Data speeds for current networks are between 2200 and 9600 bit/s using a satellite handset.
A third system was announced in 2022 when T-Mobile US and SpaceX announced a partnership to add satellite cellular service to Starlink second generation (Gen2) satellites that are to begin launching to orbit in late 2022. The service is aimed to provide dead-zone cell phone coverage across the US using existing midband PCS spectrum that T-Mobile owns. Cell coverage will begin with messaging and expand to include voice and limited data services later, with testing to begin in 2023. With Starlink Gen2 satellites in low Earth orbit using existing PCS spectrum, T-Mobile plans to be able to connect ordinary mobile phones to satellites, unlike earlier satellite phones in the market which used specialized radios to connect to geosynchronous-orbit satellites, which have longer communications latencies. T-Mobile has offered to extend the offering globally if cellular carriers in other countries wish to exchange roaming services via the T-Mobile partnership with SpaceX, with other carriers working with their regulators to enable midband communications landing rights on a country-by-country basis. Bandwidth will be limited to approximately 2 to 4 megabits per second spread across a very large cell coverage area, with thousands of voice calls or millions of text messages simultaneously in an area. The size of a single coverage area has not yet been specified.
Geotracking
LEO systems have the ability to track a mobile unit's location using Doppler navigation from the satellite. However, this method can be inaccurate by tens of kilometers. On some Iridium hardware the coordinates can be extracted using AT commands, while recent Globalstar handsets will display them on the screen.
Most VSAT terminals can be reprogrammed in-field using AT-commands to bypass automatic acquisition of GPS coordinates and instead accept manually injected GPS coordinates.
Virtual country codes
Satellite phones are usually issued with numbers in a special country calling code.
Inmarsat satellite phones are issued with codes +870. In the past, additional country codes were allocated to different satellites, but the codes +871 to +874 were phased out at the end of 2008 leaving Inmarsat users with the same country code, regardless of which satellite their terminal is registered with.
Low Earth orbit systems including some of the defunct ones have been allocated number ranges in the International Telecommunication Union's Global Mobile Satellite System virtual country code +881. Iridium satellite phones are issued with codes +881 6 and +881 7. Globalstar, although allocated +881 8 and +881 9 use U.S. telephone numbers except for service resellers located in Brazil, which use the +881 range.
Small regional satellite phone networks are allocated numbers in the +882 code designated for "international networks" which is not used exclusively for satellite phone networks.
Cost
While it is possible to obtain used handsets for the Thuraya, Iridium, and Globalstar networks for approximately , the newest handsets are quite expensive. The Iridium 9505A, released in 2001, sold in March 2010 for over $1,000. Satellite phones are purpose-built for one particular network and cannot be switched to other networks. The price of handsets varies with network performance. If a satellite phone provider encounters trouble with its network, handset prices will fall, then increase once new satellites are launched. Similarly, handset prices will increase when calling rates are reduced.
Among the most expensive satellite phones are BGAN terminals, often costing several thousand dollars. These phones provide about 0.5 Mbit/s Internet and voice communications. Satellite phones are sometimes subsidised by the provider if one signs a post-paid contract, but subsidies are usually only a few hundred dollars or less.
Since most satellite phones are built under license or the manufacturing of handsets is contracted out to OEMs, operators have a large influence over the selling price. Satellite networks operate under proprietary protocols, making it difficult for manufacturers to independently make handsets.
A startup is proposing the use of standard mobile phone technology in satellites to enable low bandwidth text message with satellites from cheap mobile phones.
Calling cost
The cost of making voice calls from a satellite phone varies from around $0.15 to $2 per minute, while calling them from landlines and regular mobile phones is more expensive. Costs for data transmissions (particularly broadband data) can be much higher. Rates from landlines and mobile phones range from $3 to $14 per minute with Iridium, Thuraya and Inmarsat being some of the most expensive networks to call. The receiver of the call pays nothing, unless they are being called via a special reverse-charge service.
Calls between different satellite phone networks are often very expensive, with calling rates of up to $15 per minute.
Calls from satellite phones to landlines are usually around $0.80 to $1.50 per minute unless special offers are used. Such promotions are usually bound to a particular geographic area where traffic is low.
Most satellite phone networks have pre-paid plans, with vouchers ranging from $100 to $5,000.
One-way services
Some satellite phone networks provide a one-way paging channel to alert users in poor coverage areas (such as indoors) of the incoming call. When the alert is received on the satellite phone it must be taken to an area with better coverage before the call can be accepted.
Globalstar provides a one-way data uplink service, typically used for asset tracking.
Iridium operates a one-way pager service as well as the call alert feature.
Legal restrictions
In some countries, possession of a satellite phone is illegal. Their signals will usually bypass local telecoms systems, hindering censorship and wiretapping attempts, which has led some intelligence agencies to believe that satellite phones aid terrorist activity. It is also common for restrictions to be in place in countries with oppressive governments regimes as a way to both expose subversive agents within their country and maximize the control of the information that makes it past their borders.
China – Inmarsat became the first company permitted to sell satellite phones in 2016. China Telecom began selling satellite phones in 2018 and six other satellite phone companies expressed their interest in entering the Chinese market shortly after.
Cuba
India – only Inmarsat-based satellite services are permitted within territories and areas under Indian jurisdiction. Importation and operation of all other satellite services, including Thuraya and Iridium, is illegal. International shipping is obliged to comply with Indian Directorate-General of Shipping (DGS) Order No. 02 of 2012 which prohibits the unauthorised import and operation of Thuraya, Iridium and other such satellite phones in Indian waters. The legislation to this effect is Section 6 of Indian Wireless Act and Section 20 of Indian Telegraph Act. International Long Distance (ILD) licences and No Objection Certificates (NOC) issued by Indian Department of Telecommunications (DOT) are mandatory for satellite communication services on Indian territory.
Mauritius – In 2022, the Information and Communications Authority started regulating the ownership and use of satellite phones.
Myanmar
North Korea – The US Bureau of Diplomatic Security advises visitors that they have "no right to privacy in North Korea and should assume your communications are monitored" which excludes the possibility of satellite phone technology.
Russia – in 2012, new regulations governing the use of satellite phones inside Russia or its territories were developed with the stated aim of fighting terrorism by enabling the Russian government to intercept calls. These regulations allow non-Russian visitors to register their SIM cards for use within Russian territory for up to six months.
Security concerns
All modern satellite phone networks encrypt voice traffic to prevent eavesdropping. In 2012, a team of academic security researchers reverse-engineered the two major proprietary encryption algorithms in use. One algorithm (used in GMR-1 phones) is a variant of the A5/2 algorithm used in GSM (used in common mobile phones), and both are vulnerable to cipher-text only attacks. The GMR-2 standard introduced a new encryption algorithm which the same research team also cryptanalysed successfully. Thus satellite phones need additional encrypting if used for high-security applications.
Use in disaster response
Most mobile telephone networks operate close to capacity during normal times, and large spikes in call volumes caused by widespread emergencies often overload the systems when they are needed most. Examples reported in the media where this has occurred include the 1999 İzmit earthquake, the September 11 attacks, the 2006 Kiholo Bay earthquake, the 2003 Northeast blackouts, Hurricane Katrina, the 2007 Minnesota bridge collapse, the 2010 Chile earthquake, and the 2010 Haiti earthquake. Reporters and journalists have also been using satellite phones to communicate and report on events in war zones such as Iraq.
Terrestrial cell antennas and networks can be damaged by natural disasters. Satellite telephony can avoid this problem and be useful during natural disasters. Satellite phone networks themselves are prone to congestion as satellites and spot beams cover a large area with relatively few voice channels.
Integration into conventional mobile phones
In the early 2020s, manufacturers began to integrate satellite connectivity into smartphone devices for use in remote areas, out of the cellular network range. The satellite-to-phone services use L band frequencies, which are compatible with most modern handsets. However, due to the antenna limitations in the conventional phones, in the early stages of implementation satellite connectivity is limited to satellite messaging and satellite emergency services.
In 2022, the Apple iPhone 14 started supporting sending emergency text messages via Globalstar satellites. In 2023, the Apple iPhone 15 added satellite communication with roadside service in the United States. In 2022, T-Mobile formed a partnership to use Starlink services via existing LTE spectrum, expected in late 2024. In 2022, AST SpaceMobile started building a 3GPP standard-based cellular space network to allow existing, unmodified smartphones to connect to satellites in areas with coverage gaps. In 2023, Qualcomm announced Snapdragon Satellite, the service that will allow supported cellphones, starting with Snapdragon 8 Gen 2 chipset, to send and receive text messages via 5G non-terrestrial networks (NTN). In 2024, Iridium introduced Project Stardust, a standard-based satellite-to-cellphone service supported via NB-IoT for 5G non-terrestrial networks, which will be utilized over Iridium's existing low-earth orbit satellites. Scheduled for launch in 2026, the service provides messaging, emergency communications and IoT for devices like cars, smartphones, tablets and related consumer applications.
| Technology | Telecommunications | null |
618077 | https://en.wikipedia.org/wiki/Power%20engineering | Power engineering | Power engineering, also called power systems engineering, is a subfield of electrical engineering that deals with the generation, transmission, distribution, and utilization of electric power, and the electrical apparatus connected to such systems. Although much of the field is concerned with the problems of three-phase AC power – the standard for large-scale power transmission and distribution across the modern world – a significant fraction of the field is concerned with the conversion between AC and DC power and the development of specialized power systems such as those used in aircraft or for electric railway networks. Power engineering draws the majority of its theoretical base from electrical engineering and mechanical engineering.
History
Pioneering years
Electricity became a subject of scientific interest in the late 17th century. Over the next two centuries a number of important discoveries were made including the incandescent light bulb and the voltaic pile. Probably the greatest discovery with respect to power engineering came from Michael Faraday who in 1831 discovered that a change in magnetic flux induces an electromotive force in a loop of wire—a principle known as electromagnetic induction that helps explain how generators and transformers work.
In 1881 two electricians built the world's first power station at Godalming in England. The station employed two waterwheels to produce an alternating current that was used to supply seven Siemens arc lamps at 250 volts and thirty-four incandescent lamps at 40 volts. However supply was intermittent and in 1882 Thomas Edison and his company, The Edison Electric Light Company, developed the first steam-powered electric power station on Pearl Street in New York City. The Pearl Street Station consisted of several generators and initially powered around 3,000 lamps for 59 customers. The power station used direct current and operated at a single voltage. Since the direct current power could not be easily transformed to the higher voltages necessary to minimise power loss during transmission, the possible distance between the generators and load was limited to around half-a-mile (800 m).
That same year in London Lucien Gaulard and John Dixon Gibbs demonstrated the first transformer suitable for use in a real power system. The practical value of Gaulard and Gibbs' transformer was demonstrated in 1884 at Turin where the transformer was used to light up forty kilometres (25 miles) of railway from a single alternating current generator. Despite the success of the system, the pair made some fundamental mistakes. Perhaps the most serious was connecting the primaries of the transformers in series so that switching one lamp on or off would affect other lamps further down the line. Following the demonstration George Westinghouse, an American entrepreneur, imported a number of the transformers along with a Siemens generator and set his engineers to experimenting with them in the hopes of improving them for use in a commercial power system.
One of Westinghouse's engineers, William Stanley, recognised the problem with connecting transformers in series as opposed to parallel and also realised that making the iron core of a transformer a fully enclosed loop would improve the voltage regulation of the secondary winding. Using this knowledge he built the world's first practical transformer based alternating current power system at Great Barrington, Massachusetts in 1886. In 1885 the Italian physicist and electrical engineer Galileo Ferraris demonstrated an induction motor and in 1887 and 1888 the Serbian-American engineer Nikola Tesla filed a range of patents related to power systems including one for a practical two-phase induction motor which Westinghouse licensed for his AC system.
By 1890 the power industry had flourished and power companies had built thousands of power systems (both direct and alternating current) in the United States and Europe – these networks were effectively dedicated to providing electric lighting. During this time a fierce rivalry in the US known as the "war of the currents" emerged between Edison and Westinghouse over which form of transmission (direct or alternating current) was superior. In 1891, Westinghouse installed the first major power system that was designed to drive an electric motor and not just provide electric lighting. The installation powered a synchronous motor at Telluride, Colorado with the motor being started by a Tesla induction motor. On the other side of the Atlantic, Oskar von Miller built a 20 kV 176 km three-phase transmission line from Lauffen am Neckar to Frankfurt am Main for the Electrical Engineering Exhibition in Frankfurt. In 1895, after a protracted decision-making process, the Adams No. 1 generating station at Niagara Falls began transmitting three-phase alternating current power to Buffalo at 11 kV. Following completion of the Niagara Falls project, new power systems increasingly chose alternating current as opposed to direct current for electrical transmission.
Twentieth century
Power engineering and Bolshevism
The generation of electricity was regarded as particularly important following the Bolshevik seizure of power. Lenin stated "Communism is Soviet power plus the electrification of the whole country." He was subsequently featured on many Soviet posters, stamps etc. presenting this view. The GOELRO plan was initiated in 1920 as the first Bolshevik experiment in industrial planning and in which Lenin became personally involved. Gleb Krzhizhanovsky was another key figure involved, having been involved in the construction of a power station in Moscow in 1910. He had also known Lenin since 1897 when they were both in the St. Petersburg chapter of the Union of Struggle for the Liberation of the Working Class.
Power engineering in the USA
In 1936 the first commercial high-voltage direct current (HVDC) line using mercury-arc valves was built between Schenectady and Mechanicville, New York. HVDC had previously been achieved by installing direct current generators in series (a system known as the Thury system) although this suffered from serious reliability issues. In 1957 Siemens demonstrated the first solid-state rectifier (solid-state rectifiers are now the standard for HVDC systems) however it was not until the early 1970s that this technology was used in commercial power systems. In 1959 Westinghouse demonstrated the first circuit breaker that used SF6 as the interrupting medium. SF6 is a far superior dielectric to air and, in recent times, its use has been extended to produce far more compact switching equipment (known as switchgear) and transformers. Many important developments also came from extending innovations in the ICT field to the power engineering field. For example, the development of computers meant load flow studies could be run more efficiently allowing for much better planning of power systems. Advances in information technology and telecommunication also allowed for much better remote control of the power system's switchgear and generators.
Power
Power Engineering deals with the generation, transmission, distribution and utilization of electricity as well as the design of a range of related devices. These include transformers, electric generators, electric motors and power electronics.
Power engineers may also work on systems that do not connect to the grid. These systems are called off-grid power systems and may be used in preference to on-grid systems for a variety of reasons. For example, in remote locations it may be cheaper for a mine to generate its own power rather than pay for connection to the grid and in most mobile applications connection to the grid is simply not practical.
Fields
Electricity generation covers the selection, design and construction of facilities that convert energy from primary forms to electric power.
Electric power transmission requires the engineering of high voltage transmission lines and substation facilities to interface to generation and distribution systems. High voltage direct current systems are one of the elements of an electric power grid.
Electric power distribution engineering covers those elements of a power system from a substation to the end customer.
Power system protection is the study of the ways an electrical power system can fail, and the methods to detect and mitigate for such failures.
In most projects, a power engineer must coordinate with many other disciplines such as civil and mechanical engineers, environmental experts, and legal and financial personnel. Major power system projects such as a large generating station may require scores of design professionals in addition to the power system engineers. At most levels of professional power system engineering practice, the engineer will require as much in the way of administrative and organizational skills as electrical engineering knowledge.
Professional societies and international standards organizations
In both the UK and the US, professional societies had long existed for civil and mechanical engineers. The Institution of Electrical Engineers (IEE) was founded in the UK in 1871, and the AIEE in the United States in 1884. These societies contributed to the exchange of electrical knowledge and the development of electrical engineering education.
On an international level, the International Electrotechnical Commission (IEC), which was founded in 1906, prepares standards for power engineering, with 20,000 electrotechnical experts from 172 countries developing global specifications based on consensus.
| Technology | Disciplines | null |
618631 | https://en.wikipedia.org/wiki/Low%20back%20pain | Low back pain | Low back pain or lumbago is a common disorder involving the muscles, nerves, and bones of the back, in between the lower edge of the ribs and the lower fold of the buttocks. Pain can vary from a dull constant ache to a sudden sharp feeling. Low back pain may be classified by duration as acute (pain lasting less than 6 weeks), sub-chronic (6 to 12 weeks), or chronic (more than 12 weeks). The condition may be further classified by the underlying cause as either mechanical, non-mechanical, or referred pain. The symptoms of low back pain usually improve within a few weeks from the time they start, with 40–90% of people recovered by six weeks.
In most episodes of low back pain a specific underlying cause is not identified or even looked for, with the pain believed to be due to mechanical problems such as muscle or joint strain. If the pain does not go away with conservative treatment or if it is accompanied by "red flags" such as unexplained weight loss, fever, or significant problems with feeling or movement, further testing may be needed to look for a serious underlying problem. In most cases, imaging tools such as X-ray computed tomography are not useful or recommended for low back pain that lasts less than 6 weeks (with no red flags) and carry their own risks. Despite this, the use of imaging in low back pain has increased. Some low back pain is caused by damaged intervertebral discs, and the straight leg raise test is useful to identify this cause. In those with chronic pain, the pain processing system may malfunction, causing large amounts of pain in response to non-serious events. Chronic non-specific low back pain (CNSLBP) is a highly prevalent musculoskeletal condition that not only affects the body, but also a person's social and economic status. It would be greatly beneficial for people with CNSLBP to be screened for genetic issues, unhealthy lifestyles and habits, and psychosocial factors on top of musculoskeletal issues. Chronic lower back pain is defined as back pain that lasts more than three months.
The symptoms of low back pain usually improve within a few weeks from the time they start, with 40–90% of people recovered by six weeks. Normal activity should be continued as much as the pain allows. Initial management with non-medication based treatments is recommended. Non–medication based treatments include superficial heat, massage, acupuncture, or spinal manipulation. If these are not sufficiently effective, NSAIDs are recommended. A number of other options are available for those who do not improve with usual treatment. Opioids may be useful if simple pain medications are not enough, but they are not generally recommended due to side effects, including high rates of addiction, accidental overdose and death. Surgery may be beneficial for those with disc-related chronic pain and disability or spinal stenosis. No clear benefit of surgery has been found for other cases of non-specific low back pain. Low back pain often affects mood, which may be improved by counseling or antidepressants. Additionally, there are many alternative medicine therapies, but there is not enough evidence to recommend them confidently. The evidence for chiropractic care and spinal manipulation is mixed.
Approximately 9–12% of people (632 million) have low back pain at any given point in time, and nearly 25% report having it at some point over any one-month period. About 40% of people have low back pain at some point in their lives, with estimates as high as 80% among people in the developed world. Low back pain is the greatest contributor to lost productivity, absenteeism, disability and early retirement worldwide. Difficulty with low back pain most often begins between 20 and 40 years of age. Women and older people have higher estimated rates of lower back pain and also higher disability estimates. Low back pain is more common among people aged between 40 and 80 years, with the overall number of individuals affected expected to increase as the population ages. According to the World Health Organizations, lower back pain is the top medical condition world-wide from which the most number of people world-wide can benefit from improved rehabilitation.
Signs and symptoms
In the common presentation of acute low back pain, pain develops after movements that involve lifting, twisting, or forward-bending. The symptoms may start soon after the movements or upon waking up the following morning. The description of the symptoms may range from tenderness at a particular point, to diffuse pain. It may or may not worsen with certain movements, such as raising a leg, or positions, such as sitting or standing. Pain radiating down the legs (known as sciatica) may be present. The first experience of acute low back pain is typically between the ages of 20 and 40. This is often a person's first reason to see a medical professional as an adult. Recurrent episodes occur in more than half of people with the repeated episodes being generally more painful than the first.
Other problems may occur along with low back pain. Chronic low back pain is associated with sleep problems, including a greater amount of time needed to fall asleep, disturbances during sleep, a shorter duration of sleep, and less satisfaction with sleep. In addition, a majority of those with chronic low back pain show symptoms of depression or anxiety.
Causes
Low back pain is not a specific disease but rather a complaint that may be caused by a large number of underlying problems of varying levels of seriousness. The majority of low back pain does not have a clear cause but is believed to be the result of non-serious muscle or skeletal issues such as sprains or strains. Obesity, smoking, weight gain during pregnancy, stress, poor physical condition, and poor sleeping position may also contribute to low back pain. There is no consensus as to whether spinal posture or certain physical activities are causal factors. A full list of possible causes includes many less common conditions. Physical causes may include osteoarthritis, degeneration of the discs between the vertebrae or a spinal disc herniation, broken vertebra(e) (such as from osteoporosis) or, rarely, an infection or tumor of the spine.
Women may have acute low back pain from medical conditions affecting the female reproductive system, including endometriosis, ovarian cysts, ovarian cancer, or uterine fibroids. Nearly half of all pregnant women report pain in the low back during pregnancy, due to changes in their posture and center of gravity causing muscle and ligament strain.
Low back pain can be broadly classified into four main categories:
Musculoskeletal – mechanical (including muscle strain, muscle spasm, or osteoarthritis); herniated nucleus pulposus, herniated disc; spinal stenosis; or compression fracture
Inflammatory – HLA-B27 associated arthritis including ankylosing spondylitis, reactive arthritis, psoriatic arthritis, inflammation within the reproductive system, and inflammatory bowel disease
Malignancy – bone metastasis from lung, breast, prostate, thyroid, among others
Infectious – osteomyelitis, abscess, urinary tract infection.
Pathophysiology
Back structures
The lumbar (or lower back) region is the area between the lower ribs and gluteal fold which includes five lumbar vertebrae (L1–L5) and the sacrum. In between these vertebrae are fibrocartilaginous discs, which act as cushions, preventing the vertebrae from rubbing together while at the same time protecting the spinal cord. Nerves come from and go to the spinal cord through specific openings between the vertebrae, receiving sensory input and sending messages to muscles. Stability of the spine is provided by the ligaments and muscles of the back and abdomen. Small joints called facet joints limit and direct the motion of the spine.
The multifidus muscles run up and down along the back of the spine, and are important for keeping the spine straight and stable during many common movements such as sitting, walking and lifting. A problem with these muscles is often found in someone with chronic low back pain, because the back pain causes the person to use the back muscles improperly in trying to avoid the pain. The problem with the multifidus muscles continues even after the pain goes away, and is probably an important reason why the pain comes back. Teaching people with chronic low back pain how to use these muscles is recommended as part of a recovery program.
An intervertebral disc has a gelatinous core surrounded by a fibrous ring. When in its normal, uninjured state, most of the disc is not served by either the circulatory or nervous systems – blood and nerves only run to the outside of the disc. Specialized cells that can survive without direct blood supply are in the inside of the disc. Over time, the discs lose flexibility and the ability to absorb physical forces. This decreased ability to handle physical forces increases stresses on other parts of the spine, causing the ligaments of the spine to thicken and bony growths to develop on the vertebrae. As a result, there is less space through which the spinal cord and nerve roots may pass. When a disc degenerates as a result of injury or disease, the makeup of a disc changes: blood vessels and nerves may grow into its interior and/or herniated disc material can push directly on a nerve root. Any of these changes may result in back pain.
Pain sensation
Pain erupts in response to a stimulus that either damages or can potentially damage the body's tissues. There are four main stages: transduction, transmission, perception, and modulation. The nerve cells that detect pain have cell bodies located in the dorsal root ganglia and fibers that transmit these signals to the spinal cord. The process of pain sensation starts when the pain-causing event triggers the endings of appropriate sensory nerve cells. This type of cell converts the event into an electrical signal by transduction. Several different types of nerve fibers carry out the transmission of the electrical signal from the transducing cell to the posterior horn of spinal cord, from there to the brain stem, and then from the brain stem to the various parts of the brain such as the thalamus and the limbic system. In the brain, the pain signals are processed and given context in the process of pain perception. Through modulation, the brain can modify the sending of further nerve impulses by decreasing or increasing the release of neurotransmitters.
Parts of the pain sensation and processing system may not function properly; creating the feeling of pain when no outside cause exists, signaling too much pain from a particular cause, or signaling pain from a normally non-painful event. Additionally, the pain modulation mechanisms may not function properly. These phenomena are involved in chronic pain.
Diagnosis
As the structure of the low back is complex, the reporting of pain is subjective, and is affected by social factors, the diagnosis of low back pain is not straightforward. While most low back pain is caused by muscle and joint problems, this cause must be separated from neurological problems, spinal tumors, fracture of the spine, and infections, among others. The ICD 10 code for low back pain is M54.5.
Classification
There are a number of ways to classify low back pain with no consensus that any one method is best. There are three general types of low back pain by cause: mechanical back pain (including nonspecific musculoskeletal strains, herniated discs, compressed nerve roots, degenerative discs or joint disease, and broken vertebra), non-mechanical back pain (tumors, inflammatory conditions such as spondyloarthritis, and infections), and referred pain from internal organs (gallbladder disease, kidney stones, kidney infections, and aortic aneurysm, among others). Mechanical or musculoskeletal problems underlie most cases (around 90% or more), and of those, most (around 75%) do not have a specific cause identified, but are thought to be due to muscle strain or injury to ligaments. Rarely, complaints of low back pain result from systemic or psychological problems, such as fibromyalgia and somatoform disorders.
Low back pain may be classified based on the signs and symptoms. Diffuse pain that does not change in response to particular movements, and is localized to the lower back without radiating beyond the buttocks, is classified as nonspecific, the most common classification. Pain that radiates down the leg below the knee, is located on one side (in the case of disc herniation), or is on both sides (in spinal stenosis), and changes in severity in response to certain positions or maneuvers is radicular, making up 7% of cases. Pain that is accompanied by red flags such as trauma, fever, a history of cancer or significant muscle weakness may indicate a more serious underlying problem and is classified as needing urgent or specialized attention.
The symptoms can also be classified by duration as acute, sub-chronic (also known as sub-acute), or chronic. The specific duration required to meet each of these is not universally agreed upon, but generally pain lasting less than six weeks is classified as acute, pain lasting six to twelve weeks is sub-chronic, and more than twelve weeks is chronic. Management and prognosis may change based on the duration of symptoms.
Red flags
The presence of certain signs, termed red flags, indicate the need for further testing to look for more serious underlying problems, which may require immediate or specific treatment. The presence of a red flag does not mean that there is a significant problem. It is only suggestive, and most people with red flags have no serious underlying problem. If no red flags are present, performing diagnostic imaging or laboratory testing in the first four weeks after the start of the symptoms has not been shown to be useful.
The usefulness of many red flags is poorly supported by evidence. The most useful for detecting a fracture are: older age, corticosteroid use, and significant trauma especially if it results in skin markings. The best determinant of the presence of cancer is a history of the same.
With other causes ruled out, people with non-specific low back pain are typically treated symptomatically, without exact determination of the cause. Efforts to uncover factors that might complicate the diagnosis, such as depression, substance abuse, or an agenda concerning insurance payments may be helpful.
Tests
Imaging is indicated when there are red flags, ongoing neurological symptoms that do not resolve, or ongoing or worsening pain. In particular, early use of imaging (either MRI or CT) is recommended for suspected cancer, infection, or cauda equina syndrome. MRI is slightly better than CT for identifying disc disease; the two technologies are equally useful for diagnosing spinal stenosis. Only a few physical diagnostic tests are helpful. The straight leg raise test is almost always positive in those with disc herniation, and lumbar provocative discography may be useful to identify a specific disc causing pain in those with chronic high levels of low back pain. Therapeutic procedures such as nerve blocks can also be used to determine a specific source of pain. Some evidence supports the use of facet joint injections, transforminal epidural injections and sacroiliac injections as diagnostic tests. Most other physical tests, such as evaluating for scoliosis, muscle weakness or wasting, and impaired reflexes, are of little use.
Complaints of low back pain are one of the most common reasons people visit doctors. For pain that has lasted only a few weeks, the pain is likely to subside on its own. Thus, if a person's medical history and physical examination do not suggest a specific disease as the cause, medical societies advise against imaging tests such as X-rays, CT scans, and MRIs. Individuals may want such tests but, unless red flags are present, they are unnecessary health care. Routine imaging increases costs, is associated with higher rates of surgery with no overall benefit, and the radiation used may be harmful to one's health. Fewer than 1% of imaging tests identify the cause of the problem. Imaging may also detect harmless abnormalities, encouraging people to request further unnecessary testing or to worry. Even so, MRI scans of the lumbar region increased by more than 300% among United States Medicare beneficiaries from 1994 to 2006.
Prevention
Exercise alone, or along with education, appears to be useful for preventing low back pain. Exercise is also probably effective in preventing recurrences in those with pain that has lasted more than six weeks. Assessing chronic low back pain, a 2007 review concluded that a firm mattress is less likely to alleviate pain compared to a medium-firm mattress, while a 2020 review stated that studies have been inadequate to comment on mattress firmness. There is little to no evidence that back belts are any more helpful in preventing low back pain than education about proper lifting techniques. Shoe insoles do not help prevent low back pain.
Studies have proven that interventions aimed to reduce pain and functional disability need to be accompanied by psychological interventions to improve a patient's motivation and attitude toward their recovery. Education about an injury and how it can effect a person's mental health is just as important as the physical rehabilitation. However, all of these interventions should occur in partnership with a structured therapeutic exercise program and assistance from a trained physical therapist.
Management
Most people with acute or subacute low back pain improve over time no matter the treatment. There is often improvement within the first month. Although fear in those suffering from low back pain often leads to avoiding activity, this is found to lead to greater disability. The recommendations include remaining active, avoiding activity that worsen the pain, and understanding self-care of the symptoms. Management of low back pain depends on which of the three general categories is the cause: mechanical problems, non-mechanical problems, or referred pain. For acute pain that is causing only mild to moderate problems, the goals are to restore normal function, return the individual to work, and minimize pain. The condition is normally not serious, resolves without much being done, and recovery is helped by attempting to return to normal activities as soon as possible within the limits of pain. Providing individuals with coping skills through reassurance of these facts is useful in speeding recovery.
For those with sub-chronic or chronic low back pain, multidisciplinary treatment programs may help. Initial management with non–medication based treatments is recommended Non–medication based treatments include superficial heat, massage, acupuncture, or spinal manipulation. If these are not sufficiently effective, NSAIDs are recommended. Acetaminophen and systemic steroids are not recommended as both medications are not effective at improving pain outcomes in acute or subacute low back pain.
Physical therapy stabilization exercises for lumbar spine and manual therapy have shown decrease in pain symptoms in patients. Manual therapy and stabilization effects have similar effects on low back pain which overweighs the effects of general exercises. The most effective types of exercise to improve low back pain symptoms are core strengthening and mixed exercise types. An appropriate type of exercise recommended is an aerobic exercise program for 12 hours of exercise over a duration of 8 weeks.
Distress due to low back pain contributes significantly to overall pain and disability experienced. Therefore, treatment strategies that aim to change beliefs and behaviours, such as cognitive-behavioural therapy can be of use.
Access to care as recommended in medical guidelines varies considerably from the care that most people with low back pain receive globally. This is due to factors such as availability, access and payment models (e.g. insurance, health-care systems).
Physical management
Management of acute low back pain
Increasing general physical activity has been recommended, but no clear relationship to pain or disability or returning to work has been found when used for the treatment of an acute episode of pain. For acute pain, low- to moderate-quality evidence supports walking. Aerobic exercises like progressive walking appears useful for subacute and acute low back pain, is strongly recommended for chronic low back pain, and is recommended after surgery. Directional exercises, which try to limit low back pain, are recommended in sub-acute, chronic and radicular low back pain. These exercises only work if they are limiting low back pain. Exercise programs that incorporate stretching only are not recommended for acute low back pain. Stretching, especially with limited range of motion, can impede future progression of treatment like limiting strength and limiting exercises. Yoga and Tai chi are not recommended in case of acute or subacute low back pain, but are recommended in case of chronic back pain.
Treatment according to McKenzie method is somewhat effective for recurrent acute low back pain, but its benefit in the short term does not appear significant. There is tentative evidence to support the use of heat therapy for acute and sub-chronic low back pain but little evidence for the use of either heat or cold therapy in chronic pain. Weak evidence suggests that back belts might decrease the number of missed workdays, but there is nothing to suggest that they help with the pain. Ultrasound and shock wave therapies do not appear effective and therefore are not recommended. Lumbar traction lacks effectiveness as an intervention for radicular low back pain. It is also unclear whether lumbar supports are an effective treatment intervention.
==== Management of chronic low back pain ====
Exercise therapy is effective in decreasing pain and improving physical function, trunk muscle strength and the mental health for those with chronic low back pain. It also improves long-term function and appears to reduce recurrence rates for as long as six months after the completion of the program. The observed treatment effect for the exercise when compared to no treatment, usual care or placebo, improved pain (low‐certainty evidence), but improvements were small for functional limitations outcomes (moderate‐certainty evidence). There is no evidence that one particular type of exercise therapy is more effective than another, so the form of exercise used can be based on patient or practitioner preference, availability and cost. The Alexander technique appears useful for chronic back pain, and there is some evidence to support small benefits from the use of yoga. If a person with chronic low back pain is motivated, it is recommended to use yoga and tai chi as a form of treatment, but this is not recommended to treat acute or subacute low back pain. Motor control exercise, which involves guided movement and use of normal muscles during simple tasks which then builds to more complex tasks, improves pain and function up to 20 weeks, but there was little difference compared to manual therapy and other forms of exercise. Motor control exercise accompanied by manual therapy also produces similar reductions in pain intensity when compared to general strength and condition exercise training, yet only the latter also improved muscle endurance and strength, whilst concurrently decreased self-reported disability. Aquatic therapy is recommended as an option in those with other preexisting conditions like extreme obesity, degenerative joint disease, or other conditions that limit progressive walking. Aquatic therapy is recommended for chronic and subacute low back pain in those with a preexisting condition. Aquatic therapy is not recommended for people that have no preexisting condition that limits their progressive walking. There is low-to-moderate quality evidence that supports pilates in low back pain for the reduction of pain and disability, however there is no conclusive evidence that pilates is better than any other form of exercise for low back pain.
Patients with chronic low back pain receiving multidisciplinary biopsychosocial rehabilitation (MBR) programs are likely to experience less pain and disability than those receiving usual care or a physical treatment. MBR also has a positive influence on work status of the patient compared to physical treatment. Effects are of a modest magnitude and should be balanced against the time and resource requirements of MBR programs.
Peripheral nerve stimulation, a minimally-invasive procedure, may be useful in cases of chronic low back pain that do not respond to other measures, although the evidence supporting it is not conclusive, and it is not effective for pain that radiates into the leg. Evidence for the use of shoe insoles as a treatment is inconclusive. Transcutaneous electrical nerve stimulation (TENS) has not been found to be effective in chronic low back pain. There has been little research that supports the use of lumbar extension machines and thus they are not recommended.
Medications
If initial management with non–medication based treatments is insufficient, medication may be recommended. As pain medications are only somewhat effective, expectations regarding their benefit may differ from reality, and this can lead to decreased satisfaction.
The medication typically prescribed first are acetaminophen (paracetamol), NSAIDs (though not aspirin), or skeletal muscle relaxants and these are enough for most people. Benefits with NSAIDs is thought to be small, but is more effective than Acetaminophen (paracetamol), which may be no more effective than placebo at improving pain, quality of life, or function. NSAIDs however, carry a greater risk of side effects, including kidney failure, stomach ulcers and possibly heart problems, so it is used at the lowest effective dosage for the shortest possible time. NSAIDs are available in several different classes; there is no evidence to support the use of COX-2 inhibitors over any other class of NSAIDs with respect to benefits. With respect to safety naproxen may be best. Muscle relaxants may be beneficial.
Systemic corticosteriods are sometimes suggested for low back pain and may have a small benefit in the short-term for radicular low back pain, however, the benefit for non-radicular back pain and the optimal dose and length of treatment is unclear.
As of 2022, the CDC has released a guideline for prescribed opioid use in the management of chronic pain. It states that opioid use is not the preferred treatment when managing chronic pain due to the excessive risks involved, including high risks of addiction, accidental overdose and death. Specialist groups advise against general long-term use of opioids for chronic low back pain. If the pain is not managed adequately, short-term use of opioids such as morphine may be suggested, although low back pain outcomes are poorer in the long-term. If prescribed, a person and their clinician should have a realistic plan to discontinue its use in the event that the risks outweigh the benefit. These medications carry a risk of addiction, may have negative interactions with other drugs, and have a greater risk of side effects, including dizziness, nausea, and constipation. Opioid treatment for chronic low back pain increases the risk for lifetime illicit drug use and the effect of long-term use of opioids for lower back pain is unknown. For older people with chronic pain, opioids may be used in those for whom NSAIDs present too great a risk, including those with diabetes, stomach or heart problems. They may also be useful for a select group of people with neuropathic pain.
Antidepressants may be effective for treating chronic pain associated with symptoms of depression, but they have a risk of side effects. Although the antiseizure drugs gabapentin, pregabalin, and topiramate are sometimes used for chronic low back pain evidence does not support a benefit. Systemic oral steroids have not been shown to be useful in low back pain. Facet joint injections and steroid injections into the discs have not been found to be effective in those with persistent, non-radiating pain; however, they may be considered for those with persistent sciatic pain. Epidural corticosteroid injections provide a slight and questionable short-term improvement in those with sciatica but are of no long-term benefit. There are also concerns of potential side effects.
Surgery
Surgery may be useful in those with a herniated disc that is causing significant pain radiating into the leg, significant leg weakness, bladder problems, or loss of bowel control. It may also be useful in those with spinal stenosis. In the absence of these issues, there is no clear evidence of a benefit from surgery.
Discectomy (the partial removal of a disc that is causing leg pain) can provide pain relief sooner than nonsurgical treatments. Discectomy has better outcomes at one year but not at four to ten years. The less invasive microdiscectomy has not been shown to result in a different outcome than regular discectomy. For most other conditions, there is not enough evidence to provide recommendations for surgical options. The long-term effect surgery has on degenerative disc disease is not clear. Less invasive surgical options have improved recovery times, but evidence regarding effectiveness is insufficient.
For those with pain localized to the lower back due to disc degeneration, fair evidence supports spinal fusion as equal to intensive physical therapy and slightly better than low-intensity nonsurgical measures. Fusion may be considered for those with low back pain from acquired displaced vertebra that does not improve with conservative treatment, although only a few of those who have spinal fusion experience good results, and there may be no clinically important difference between disk replacement and fusion surgery. There are a number of different surgical procedures to achieve fusion, with no clear evidence of one being better than the others. Adding spinal implant devices during fusion increases the risk but provides no added improvement in pain or function. Spinal cord stimulation using implanted electrodes is not supported by evidence due to the potential risks and costs.
Alternative medicine
It is unclear if alternative treatments are useful for non-chronic back pain. Chiropractic care or spinal manipulation therapy (SMT) appear similarly effective to other recommended treatments. National guidelines differ, with some not recommending SMT, some describing manipulation as optional, and others recommending a short course for those who do not improve with other treatments. A 2017 review recommended SMT based on low-quality evidence. There is insufficient evidence to recommend manipulation under anaesthesia, or medically assisted manipulation. SMT does not provide significant benefits compared to motor control exercises.
The evidence supporting acupuncture treatment for providing clinically beneficial acute and chronic pain relief is very weak. When compared to a 'sham' treatment, no differences in pain relief or improvements in a person's quality of life were found. There is very weak evidence that acupuncture may be better than no treatment at all for immediate relief. A 2012 systematic review reported the findings that for people with chronic pain, acupuncture may improve pain a little more than no treatment and about the same as medications, but it does not help with disability. This pain benefit is only present right after treatment and not at follow-up. Acupuncture may be an option for those with chronic pain that does not respond to other treatments like conservative care and medications, however this depends on patient preference, the cost, and on how accessible acupuncture is for the person.
Massage therapy does not appear to provide much benefit for acute low back pain. Massage therapy has been found to be more effective for acute low back pain than no treatment; the benefits were found to be limited to the short term and there was no effect for improving function. For chronic low back pain, massage therapy was no better than no treatment for both pain and function, though only in the short-term. The overall quality of the evidence was low and the authors conclude that massage therapy is generally not an effective treatment for low back pain. Massage therapy is recommended for selected people with subacute and chronic low back pain, but it should be paired with another form of treatment like aerobic or strength exercises. For acute or chronic radicular pain syndromes massage therapy is recommended only if low back pain is considered a symptom. Mechanical massage tools are not recommended for the treatment of any form of low back pain.
Prolotherapy – the practice of injecting solutions into joints (or other areas) to cause inflammation and thereby stimulate the body's healing response – has not been found to be effective by itself, although it may be helpful when added to another therapy.
Herbal medicines, as a whole, are poorly supported by evidence. The herbal treatments Devil's claw and white willow may reduce the number of individuals reporting high levels of pain; however, for those taking pain relievers, this difference is not significant. Capsicum, in the form of either a gel or a plaster cast, has been found to reduce pain and increase function.
Behavioral therapy may be useful for chronic pain. There are several types available, including operant conditioning, which uses reinforcement to reduce undesirable behaviors and increase desirable behaviors; cognitive behavioral therapy, which helps people identify and correct negative thinking and behavior; and respondent conditioning, which can modify an individual's physiological response to pain. The benefit however is small. Medical providers may develop an integrated program of behavioral therapies. The evidence is inconclusive as to whether mindfulness-based stress reduction reduces chronic back pain intensity or associated disability, although it suggests that it may be useful in improving the acceptance of existing pain.
Tentative evidence supports neuroreflexotherapy (NRT), in which small pieces of metal are placed just under the skin of the ear and back, for non-specific low back pain. Multidisciplinary biopsychosocial rehabilitation (MBR), targeting physical and psychological aspects, may improve back pain but evidence is limited. There is a lack of good quality evidence to support the use of radiofrequency denervation for pain relief.
KT Tape has been found to be no different for management of chronic non-specific low back pain than other established pain management strategies.
Education
There is strong evidence that education may improve low back pain, with a 2.5 hour educational session more effective than usual care for helping people return to work in the short- and long-term. This was more effective for people with acute rather than chronic back pain. The benefit of training for preventing back pain in people who work manually with materials or is not clear, however moderate quality evidence does not show a role in preventing back pain.
Prognosis
Overall, the outcome for acute low back pain is positive. Pain and disability usually improve a great deal in the first six weeks, with complete recovery reported by 40 to 90%. In those who still have symptoms after six weeks, improvement is generally slower with only small gains up to one year. At one year, pain and disability levels are low to minimal in most people. Distress, previous low back pain, and job satisfaction are predictors of long-term outcome after an episode of acute pain. Certain psychological problems such as depression, or unhappiness due to loss of employment may prolong the episode of low back pain. Following a first episode of back pain, recurrences occur in more than half of people.
For persistent low back pain, the short-term outcome is also positive, with improvement in the first six weeks but very little improvement after that. At one year, those with chronic low back pain usually continue to have moderate pain and disability. People at higher risk of long-term disability include those with poor coping skills or with fear of activity (2.5 times more likely to have poor outcomes at one year), those with a poor ability to cope with pain, functional impairments, poor general health, or a significant psychiatric or psychological component to the pain (Waddell's signs).
Prognosis may be influenced by expectations, with those having positive expectations of recovery related to higher likelihood of returning to work and overall outcomes.
Epidemiology
Low back pain that lasts at least one day and limits activity is a common complaint. Globally, about 40% of people have low back pain at some point in their lives, with estimates as high as 80% of people in the developed world. Approximately 9 to 12% of people (632 million) have low back pain at any given point in time, which was calculated to 7460 per 100,000 globally in 2020. Nearly one quarter (23.2%) report having it at some point over any one-month period. Difficulty most often begins between 20 and 40 years of age. However, low back pain becomes increasingly common with age, and is most common in the age group of 85. Older adults more greatly affected by low back pain; they are more likely to lose mobility and independence and less likely to continue to participate in social and family activities.
Women have higher rates of low back pain than men within all age groups, and this difference becomes more marked in older age groups (above 75 years). In a 2012 review which found a higher rate in females than males, the reviewers thought this may be attributable to greater rates of pains due to osteoporosis, menstruation, and pregnancy among women, or possibly because women were more willing to report pain than men. An estimated 70% of women experience back pain during pregnancy with the rate being higher the further along in pregnancy.
Although the majority of low back pain has no specific underlying cause, workplace ergonomics, smoking and obesity are associated with low back pain in approximately 30% of cases. Low levels of activity is also associated with low back pain. Workplace ergonomics associated with low back pain include lifting, bending, vibration and physically demanding work, as well as prolonged sitting, standing and awkward postures. Current smokers – and especially those who are adolescents – are more likely to have low back pain than former smokers, and former smokers are more likely to have low back pain than those who have never smoked.
The overall number of individuals affected expected to increase with population growth and as the population ages, with the largest increases expectedin low- and middle-income countries.
History
Low back pain has been with humans since at least the Bronze Age. The oldest known surgical treatise – the Edwin Smith Papyrus, dating to about 1500 BCE – describes a diagnostic test and treatment for a vertebral sprain. Hippocrates ( – ) was the first to use a term for sciatic pain and low back pain; Galen (active mid to late second century CE) described the concept in some detail. Physicians through the end of the first millennium recommended watchful waiting. Through the Medieval period, folk medicine practitioners provided treatments for back pain based on the belief that it was caused by spirits.
At the start of the 20th century, physicians thought low back pain was caused by inflammation of or damage to the nerves, with neuralgia and neuritis frequently mentioned by them in the medical literature of the time. The popularity of such proposed causes decreased during the 20th century. In the early 20th century, American neurosurgeon Harvey Williams Cushing increased the acceptance of surgical treatments for low back pain. In the 1920s and 1930s, new theories of the cause arose, with physicians proposing a combination of nervous system and psychological disorders such as nerve weakness (neurasthenia) and female hysteria. Muscular rheumatism (now called fibromyalgia) was also cited with increasing frequency.
Emerging technologies such as X-rays gave physicians new diagnostic tools, revealing the intervertebral disc as a source for back pain in some cases. In 1938, orthopedic surgeon Joseph S. Barr reported on cases of disc-related sciatica improved or cured with back surgery. As a result of this work, in the 1940s, the vertebral disc model of low back pain took over, dominating the literature through the 1980s, aiding further by the rise of new imaging technologies such as CT and MRI. The discussion subsided as research showed disc problems to be a relatively uncommon cause of the pain. Since then, physicians have come to realize that it is unlikely that a specific cause for low back pain can be identified in many cases and question the need to find one at all as most of the time symptoms resolve within 6 to 12 weeks regardless of treatment.
Society and culture
Low back pain results in large economic costs. In the United States, it is the most common type of pain in adults, responsible for a large number of missed work days, and is the most common musculoskeletal complaint seen in the emergency department. In 1998, it was estimated to be responsible for $90 billion in annual health care costs, with 5% of individuals incurring most (75%) of the costs. Between 1990 and 2001 there was a more than twofold increase in spinal fusion surgeries in the US, despite the fact that there were no changes to the indications for surgery or new evidence of greater usefulness. Further costs occur in the form of lost income and productivity, with low back pain responsible for 40% of all missed work days in the United States. Low back pain causes disability in a larger percentage of the workforce in Canada, Great Britain, the Netherlands and Sweden than in the US or Germany. In the United States, low back pain is highest of Years Lived With Disability (YLDs) rank, rate, and rercentage Change for the 25 leading causes of disability and injury, between 1990 and 2016.
Workers who experience acute low back pain as a result of a work injury may be asked by their employers to have x-rays. As in other cases, testing is not indicated unless red flags are present. An employer's concern about legal liability is not a medical indication and should not be used to justify medical testing when it is not indicated. There should be no legal reason for encouraging people to have tests which a health care provider determines are not indicated.
Research
Total disc replacement is an experimental option, but no significant evidence supports its use over lumbar fusion. Researchers are investigating the possibility of growing new intervertebral structures through the use of injected human growth factors, implanted substances, cell therapy, and tissue engineering.
| Biology and health sciences | Disabilities | Health |
618854 | https://en.wikipedia.org/wiki/Amsterdam%20Metro | Amsterdam Metro | The Amsterdam Metro () is a rapid transit system serving Amsterdam, Netherlands, and extending to the surrounding municipalities of Diemen and Ouder-Amstel. Until 2019, it also served the municipality of Amstelveen, but this route was closed and converted into a tram line. The network is owned by the City of Amsterdam and operated by municipal public transport company Gemeente Vervoerbedrijf (GVB), which also operates trams, free ferries and local buses.
The metro system consists of five routes and serves 39 stations, with a total length of . Three routes start at Amsterdam Centraal: Route 53 and Route 54 connect the city centre with the suburban residential towns of Diemen, Duivendrecht and Amsterdam-Zuidoost (the city's southeastern borough), while Route 51 first runs south and then follows a circular route connecting the southern and western boroughs. Route 50 connects Zuidoost to the Amsterdam-West borough using a circular line, which it shares with Route 51. It is the only route that does not cross the city centre. A fifth route, Route 52, running from the Amsterdam-Noord (north) borough to Amsterdam-Zuid (south) via Amsterdam Centraal, came into operation on 21 July 2018. As opposed to the other routes, it runs mostly through bored tunnels and does not share tracks with any other route.
History
Planning history
The first plans for an underground railway in Amsterdam date from the 1920s: in November 1922, members of the municipal council of Amsterdam Zeeger Gulden and Emanuel Boekman asked the responsible alderman Ter Haar to study the possibility of constructing an underground railway in the city, in response to which the municipal department of Public Works drafted reports with proposals for underground railways in both 1923 and 1929. These plans stalled in the planning phase, however, and it took until the 1950s for the discussion about underground rail to resurface again in Amsterdam.
The post-war population boom and increase in motorized traffic shifted the perception of underground rail transport in Amsterdam considerably: whereas in the 1920s, underground rail had been considered too expensive, halfway through the 1950s it was presented as a realistic solution to the problems caused by increased traffic. In 1955, a report published by the municipal government concerning the inner city of Amsterdam—known by the Dutch title Nota Binnenstad—suggested installing a commission to explore solutions to the traffic problems Amsterdam faced. This commission, which was headed by former director of the department of Public Works J.W. Clerx, was subsequently installed in March 1956, and published its report Openbaar vervoer in de agglomeratie Amsterdam in 1960.
The aldermen and mayor of Amsterdam agreed with the conclusion of the report of the Clerx commission that an underground railway network ought to be built in Amsterdam in the near future. In April 1963 they installed the Bureau Stadsspoorweg which had the task to study the technical feasibility of a metropolitan railway, to propose a route network, to suggest the preferred order of construction of the various lines, and to study the adverse effects of constructing a metro line, such as traffic disruption and the demolition of buildings.
In 1964 and 1965, Bureau Stadsspoorweg presented four reports to the municipal government of Amsterdam, which were made available to the public on 30 August 1966. In March 1968, the aldermen and mayor of Amsterdam subsequently submitted a proposal to the municipal council of Amsterdam to agree to the construction of the metro network, which the council assented to on 16 May 1968 with 38 votes in favour and 3 against. Under the original plan, four lines were to be built, connecting the entire city and replacing many of the existing tram lines. The following lines were planned: an east–west line from the southeast to the Osdorp district via Amsterdam Centraal railway station; a circle line from the western harbor area to the suburban town of Diemen; a north–south line from the northern district via Amsterdam Centraal to Weteringplantsoen traffic circle, with two branches at both ends; and a second east–west line from Geuzenveld district to Gaasperplas. The system would be constructed gradually and was expected to be completed by the end of the 1990s.
Design and construction
The first part of the original plan to be carried out was the construction of the Oostlijn (East Line), which started in 1970. The East Line links the city centre with the large-scale Bijlmermeer residential developments in the south-east of the city. It opened in 1977. The East Line starts underground, crossing the city centre and adjacent neighbourhoods in the eastern districts until Amsterdam Amstel railway station, where it continues above ground in southeastern direction. At Van der Madeweg metro station, the line splits into two branches: the Gein Branch for Route 54 and Gaasperplas Branch for Route 53. Since 1980, the northern terminus for both routes is Amsterdam Centraal railway station.
Ben Spangberg, an architect from the Government of Amsterdam, was asked to design the stations on the line. After two years of work, he told that it was too much for one person, and was assigned to the project as well. Spangberg said that they avoided sharp corners and used smooth designs instead. He wanted to have elevators in all stations, but initially wasn't allowed due to cost issues and because Nederlandse Spoorwegen didn't want to stay behind. After further discussions, the architects were permitted to design two elevator shafts, with only one of them active. Construction on the tunnels started while the stations were still being designed. The architects frequently visited the construction sites to instruct the workers to do something specific to allow for possible changes in the future. Most underground areas of the line were built by using long caissons. The caissons were built on the spot and lowered into their place by blasting away the soil beneath. This method required the demolition the houses above the tunnel. The Wibautstraat station and the tunnels near it were constructed differently, using the cut-and-cover method.
During the construction of the metro tunnel, the decision to demolish the Nieuwmarkt neighbourhood in the city centre led to strong protests in the spring of 1975 from action groups consisting of locals and members of the highly active Amsterdam squatting movement. Wall decorations at the Nieuwmarkt metro station are a reference to the protests, which are known as the Nieuwmarkt riots (Nieuwmarktrellen). Despite the protests, construction of the metro line continued but plans to build a highway through the area were abandoned. In addition, the original plans for an east–west metro line were cancelled. One of the sites where this line was to connect with the East Line had already been built underneath Weesperplein station. This lower level of Weesperplein station was never opened to the public, but its existence can still be noticed by the elevator buttons. Since the East Line was planned and built during the Cold War, Weesperplein station also features a bomb shelter which has never been used as such.
Later lines
In 1990, the Amstelveenlijn (Amstelveen Line) was opened, which is used for Route 51. Under a political compromise between the city of Amsterdam and the municipality of Amstelveen, the northern section of the line was built as a metro line while the southern section is an extended tram line. Therefore, Route 51 was originally referred to as a 'sneltram' (express tram) service, and the vehicles were manufactured to light rail standards. The changeover between third rail and overhead tramline power took place at Zuid Station.
From March 2019 onwards, the Amstelveenlijn as a sneltram ceased to exist, and is being replaced by a tram line terminating at Zuid station, including a €300 million rebuild of the original line. For connection to the metro, passengers will have to transfer at Zuid station. Line number 51 was retained for a new circular line between Isolatorweg and Central Station.
In 1997, the Ringlijn (Ring Line), which is used for Route 50, was added to the system. The line provides a rapid transit connection between the south and the west of the city, eliminating the necessity of crossing the city centre.
In 2018, the Noord-Zuidlijn (North–South Line) was added to the network. The line provides a fast connection from the north of the IJ waterway to the south of Amsterdam.
Network
From 1997 to 2018 the Amsterdam metro system consisted of four metro routes. The oldest routes are Route 54 (from Centraal station to Gein) and Route 53 (from Centraal station to Gaasperplas). Both routes are using the Oostlijn (East Line) infrastructure, which was completed in 1977. Route 51 (from Centraal station to Amstelveen Westwijk), using part of the East Line as well as the Amstelveenlijn (Amstelveen Line), was added in 1990. Route 50 (from Isolatorweg to Gein) using the Ringlijn (Ring Line or Circle Line), which was completed in 1997, as well as part of the East Line infrastructure.
A fifth line, Route 52 (from Noord station to Zuid station), was added to the network operating the Noord-Zuidlijn (North–South Line), which was completed and opened on 21 July 2018.
There are 33 full metro stations, Since Route 52 on the new North-South Line opened, six additional stations and of route have been added to the metro system, yielding a new combined network length of . In 2019, sneltram Route 51 no longer operates into the metro network. The southern sneltram portion was closed for conversion to be incorporated into the tram network.
East Line (Routes 53 and 54)
Route
On 14 October 1977, the first metro train ran on the Oostlijn (East Line) from Weesperplein to Amsterdam-Zuidoost, with two branches respectively going to Gaasperplas (now Route 53) and Holendrecht (now Route 54). Spaklerweg station was completed as a shell, but opened later. On 11 October 1980, both routes were extended to Amsterdam Centraal Station, which is now their northern terminus. The Gein Branch was extended in the southern direction on 27 August 1982, when the section between Holendrecht and Gein was completed. Spaklerweg station was then opened. In some plans for the Gein Branch, an extension to Weesp and Almere was being considered. According to the most recent regional planning study, that now seems unlikely.
Architecture
A notable part of the East Line infrastructure is a dual metro overpass on the Gaasperplas Branch in the Bijlmermeer district between Ganzenhoef and Kraaiennest stations. This long colonnade contains two single crossovers, each consisting of 33 pillars carrying long beams. The center-to-center distance between the two overpasses is . This exceptional height was necessary because the metro had to bridge the main thoroughfares in the Bijlmermeer district which were built on a system of raised embankments and viaducts.
The stations, the infrastructure and the Diemen-Zuid maintenance facility of the East Line were designed by Ben Spängberg and Sier van Rhijn, two architects at the former Public Works Department of the City of Amsterdam. Their designs in Brutalist style are characterized by large-scale application of bare concrete and excessive space in the underground station halls. It also included a sophisticated use of colour. For example, the red colour of the train doors in the original design was also used at major facilities such as billboards, gates, elevator doors, bins, and the platform signage. For the design of the entire East Line Spängberg and Van Rhijn received Merkelbach Award in 1979. The East Line was also awarded the Betonprijs (Concrete Award) in 1981, which is commemorated by the award plaques in the concourses of Centraal station and Gein station.
As part of the city's policy that one percent of construction budgets for public works had to be spent on art, all stations on the Oostlijn have been decorated by different artists. In addition, the western wall of the tunnel was painted with lines and patterns which altered between the two stations, providing passengers with a fascinating view during the ride. Over the years, these decorations have completely been covered with graffiti. Some of the station artworks have also disappeared. Plans to remove all artworks as part of the large-scale renovation of the East Line tunnel in 2012 were altered after citizens' protests stating their historical significance.
Renovation
Over the years, several stations along the East Line were expanded or renovated. Since 2003, the metro station at Amsterdam Centraal station has been continuously under construction in order to accommodate the new North–South Line station. As part of commercial development of the area surrounding the Amsterdam Arena football stadium, which included a new major business and shopping district, the Bijlmer Arena station was substantially enlarged in order to handle the increasing number of passengers. The new station, designed by Grimshaw and Arcadis Articon Architects, opened in 2007 and was shortlisted for the Stirling Prize of the Royal Institute of British Architects. Another East Line station, Kraaiennest on Route 53, was reconstructed and upgraded in 2013 as part of the major urban renewal efforts in the Bijlmermeer district. The station designed by Maccreanor Lavington features a stainless steel facade with a floral design, which, according to the architects, "allows the station to be a lantern for the local neighbourhood, creating a sense of warmth on street level and creating an instantly recognisable feature for the station" at night time. The design was awarded the 2014 EU Stirling Prize.
A major overhaul of sixteen East Line stations was announced in June 2014. The renovation works taking place from 2015 until 2017 should bring more light and space to the stations. By removing paint layers on the walls, the original Brutalist architecture will become more pronounced. In addition, disused ticket offices are to be removed and lighting and signage will be improved.
Amstelveen Line (former route 51)
History
Following the Nieuwmarkt Riots in 1975, the next major expansion of the metro network into the bordering city of Amstelveen was politically sensitive. When the decision was made to begin construction of the Amstelveenlijn (Amstelveen Line) in 1984, it was originally considered an express tram service rather than a fully-fledged metro route. On 1 December 1990 the section running from Spaklerweg to Poortwachter Station in Amstelveen was completed. As the sensitivity surrounding the metro expansion waned in the 1990s, the route was increasingly being referred to as a metro service. On 13 September 2004 an extension to Amstelveen Westwijk was completed.
Originally, the entire route of the Amstelveen Line from Spaklerweg, where it connects with the East Line, to Amstelveen was to be powered via overhead wiring. Eventually, it was decided to use a third rail between Spaklerweg and Zuid station in order to be able to increase metro service on this section of the line during exhibitions at the RAI convention centre, and overhead wiring on the southern section into Amstelveen.
The line was officially opened on 30 November 1990, replacing the overcrowded bus route 67. The equipment, lightrail series S1 and S2, was built between 1990 and 1994 by Belgian manufacturer BN in Bruges. From 1994, a total of 25 light rail vehicles was in operation. Since the extension to Westwijk in 2004, a number of S3 series trains are sometimes used on this route, raising the total number of vehicles available to 29.
Shortly before the opening, two lightrail vehicles had collided during trial runs, which reduced the number of vehicles available for the route to 11. Because of the lack of equipment and startup problems, the route was initially operated with limited service. In February 1991 heavy snowfall, continuing technical problems and equipment shortages led to the decision to limit the route to the route between Centraal Station and Zuid Station for nearly seven months, with replacement bus services on the remaining route into Amstelveen.
In 2018, after the completion of the Noord-Zuidlijn (North–South Line), there would be no more room at Amsterdam Zuid station for Route 51 to continue as express tram service into Amstelveen. According to a long-term regional planning study of 2011, the Amstelveenlijn was to be upgraded to a fully-fledged metro service. On 12 March 2013, however, the regional council of the City Region of Amsterdam decided that Route 51 would be replaced by an improved express tram service running from Westwijk to Zuid Station and a separate metro service running from Zuid Station and Amstel Station. Passengers from Amstelveen would then be required to change at Zuid Station to the metro route for Amstel or to the new Route 52 for Centraal Station. It was also decided that Tram Route 5 would run between Amstelveen and Westergasfabriek. Conversion of the southern section of the Amstelveen Line to tram operation started in 2016.
On 3 March 2019, the Amstelveen branch (the hybrid metro/tram line) was cut from route 51 as the tunnel connecting the metro line with the tram network had to make way for an underground section of the A10 motorway. In December 2020, the Amstelveen Line would become tram line 25.
Route
From Centraal Station to Amsterdam Zuid station, Route 51 ran as a full metro service and had no at-grade intersections. The light-rail vehicles on the line were powered by a third rail with the line being suitable for wide trams. The BN vehicles, however, had a width of which was the maximum width on the southern section of the line between Zuid Station and Westwijk. In order to bridge the gap between the trains and the platforms in northern section, the vehicles were equipped with retractable footboards at the doors. In addition, the vehicles were equipped with pantographs in order to retrieve power from the overhead wiring on the southern section. From Zuid Station to Westwijk, the route operated as an express tram service. On the northern half of this section, Route 51 shared tracks with Line 5 of the tramway, with dual height platforms provided at the stops shared by both lines.
Since 3 March 2019, the Amstelveen section of route 51 has been discontinued. Line 51 was upgraded to a full metro line and now runs between Central Station and Isolatorweg. Since December 2020, tram line 25 has been serving the route from Zuid Station to Westwijk.
Ring Line (Route 50)
Opened on 1 June 1997, the Ringlijn (Ring Line or Circle Line) is entirely built on embankments and viaducts, and has no level crossings. The line was initially for political reasons called "express circle tram", but since the opening of the Ring Line the transit service on the line is referred to as a Metro Route 50 (from Gein to Isolatorweg). Because it was originally considered a tram line, the light rail vehicle width of 2,65 meters was to be applied; the width that was also used on the Amstelveen Line. The new "trams" (Series M4 and S3) have retractable running boards to bridge the space between the vehicle and the platform at existing stations. Since Route 50 proved hugely popular, the express tram vehicles were insufficient to handle the number of passengers. Instead of ordering additional vehicles, in 2000 the city of Amsterdam decided to adjust the platforms at the stations between Amstelveenseweg and Isolatorweg, whereby the older rolling stock (M1, M2 and M3) serving on the East Line could also serve on the Ring Line. Such an operation was already taken into account during the construction of the stations.
North–South line (Route 52)
In 2002, the construction of the Noord/Zuidlijn (North–South line) was started. The new metro line is the first to serve the Amsterdam North district, via a tunnel under the IJ. From there, it runs via Amsterdam Centraal to Amsterdam Zuid, which is planned to become the second biggest transport hub in the city, after Amsterdam Centraal. The line includes a mixture of bored tunnels and immersed tunnels under the IJ.
The construction programme experienced several difficulties, mainly at Amsterdam Centraal, resulting in the project running more than 40% over budget and the opening being delayed several times. The project initially had a budget of €1.46 billion, but after several setbacks the total cost estimate has been adjusted to €3.1 billion (at 2009 prices). The original planned opening was for 2011, but eventually the line was opened on 21 July 2018.
The North–South line might be extended to Amsterdam Airport Schiphol in the future. In August 2014, it was announced that the line was to be equipped with 4G mobile phone coverage, to be funded jointly by the major mobile phone operators.
Planned expansion
The tram line to IJburg in the east was originally planned to be a metro line. For this reason, a short tunnel was constructed eastwards from Centraal Station underneath the railway lines. As this line was eventually constructed for tram services, the tunnel was abandoned, and there are plans to use it as part of a chocolate museum. There are still plans for the tram to IJburg to be upgraded to metro and connect to the nearby city of Almere.
On completion of the north–south metro line, Amsterdam Municipality announced it was analysing a possible east–west line at a projected cost of €7 billion.
In January 2019, CEO of Amsterdam Airport Schiphol Dick Benschop announced that agreements had been reached to extend the north–south line to the airport and Hoofddorp. In June 2019, the province of North Holland outlined plans to extend the metro to Zaandam and Purmerend along with Schiphol and Hoofddorp. A station box has already been constructed for a potential underground station in Sixhaven, on the north–south line between Noorderpark and Centraal, to be opened at a later date.
Technology
Rolling stock
As of January 2016, there are 90 electric multiple unit train sets in use within the Amsterdam Metro system. All use standard gauge track and operate on a 750 V DC third rail electrification system.
M1/M2/M3
The original, first-generation fleet consisted of types M1, M2 and M3, designed as four-axle, two-car sets manufactured by the German firm Linke-Hofmann-Busch and delivered between 1973 and 1980. These first-generation trains are nicknamed zilvermeeuw (herring gull) because of their body of unpainted steel creating a silvery look. In 2009, all trains were provided with a new interior design by different artists. As they are built to full metro carrying capacity, they were used mainly on the east line services, Routes 53 and 54, with occasional use on Route 50.
As they neared the end of their life cycle and spare parts no longer became available, the entire fleet of M1-M3 trains was gradually taken out of service permanently from 2012 to 2015, being replaced by the modern M5 trains. The last unit (No. 23) was retired after a farewell tour on 19 December 2015 and has been preserved as a heritage train. All the other units have been scrapped, with the last of these being scrapped in December 2015.
S1/S2/S3/M4
Until the arrival of the new M5 units, the remainder of the fleet consisted of smaller, narrower two section, 6-axle units that could operate both on the main metro network and the light rail ("sneltram") line to Amstelveen. Types S1 and S2, manufactured by La Brugeoise et Nivelles in Belgium, were the first units to be produced for use on this new line. In service since 1990, they currently operate exclusively on Route 51, although they could technically also be used on other lines though this has never been done. These vehicles are equipped with both third rail shoes and pantographs, along with retractable footboards to bridge the gap between the trains and the platforms on sections built to full metro standards. They are due to be withdrawn by 2020 with the conversion of the Amstelveen line to express tram service.
Types M4 and S3 were built by CAF in Spain to expand the fleet and have been in service since 1997. Type M4 was built for the new Ring Line service to Isolatorweg and is hence only third rail equipped. They mainly operate on Route 50 but can also be found on Routes 53 and 54, but never on route 51 due to their lack of pantographs. Four vehicles of the same design, designated type S3, have been equipped with pantographs to also serve on Route 51, but these rarely appear on other lines. As the platforms on the Ring Line were originally built to a smaller loading gauge than those on Routes 53 and 54, M4 and S3 sets were also equipped with retractable footboards to permit boarding on the section that Route 50 shares with Route 54. When the platforms on the Ring Line were narrowed to accommodate the older but wider M1-M3 sets, the boards were permanently removed on all M4 sets, but not on the S3 sets due to the limited loading gauge of the Amstelveen line.
M5
A newer addition to the Amsterdam Metro fleet is the M5 series, manufactured by the Polish manufacturer Alstom Konstal based on its Metropolis family of high-capacity metro trains, variants of which are in use in several foreign metro systems. Delivered from June 2013 onwards, the M5 series departs radically from previous generation units by coming in six-car articulated sets with gangway connections between all cars. Although the trains are suitable for unmanned service, they remain controlled by drivers for the time being. However, the trains are compatible with Alstom's "Urbalis" communications-based train control system which will replace the current signalling system by 2017 and enable automatic train operation across the entire network. The M4 sets have been similarly equipped with this system in early 2016.
The initial order of 28 M5 metro sets, each carrying up to 1,000 passengers, was placed to replace all M1-M3 sets on the East Line as well as to increase overall capacity on the generally overstretched metro network. As such, they are used on all routes except Route 51. For Route 52 on the North–South Line an option was taken on a second series of 12 trains which was originally designated M6. However, the GVB now refers to all trains of this type as the M5 series.
M7
The newest addition to the Amsterdam Metro fleet is the M7 series, manufactured by the Spanish manufacturer CAF, which entered passenger service on February 28, 2023. Onwards, the M7 series departs radically from previous generation M5 units by coming in three-car articulated sets with gangway connections between all cars.
The initial order of 30 M7 metro sets, was placed to replace all S1-S2 in 2024, and S3-M4 in 2027. As such, they are used on line 50, and later on 53 and 54.
Summary
Ticketing system
The OV-chipkaart, a nationwide contactless smart card system, is the only valid form of ticket on the metro system. It replaced the so-called strippenkaart system on 27 August 2009, after the two systems had run parallel since 2006. Ticket barriers have been installed in all metro stations, with free-standing card readers where platforms are shared with train or tram lines. Amstelveen Line light rail stations are only equipped with free-standing card readers.
Graphic design
Signage on the Amsterdam metro system has featured multiple designs stemming from different eras. The original 1974 signage uses the M.O.L. typeface, which was specially designed for the metro by Gerard Unger. The openings within the letters are larger than normal in order to improve the letters' legibility when illuminated. The name M.O.L. refers to the Dutch word mol which means mole in English. The idea to use a mole as the mascot for the metro was rejected by city authorities. Other versions are the 1991 version found on the Amstelveen Line, the 1995 version found mainly on Ring Line and the 2009 version which has replaced earlier versions at many stations. In 2016 the City Region of Amsterdam commissioned a new signage system and logo in an effort to harmonize all of the signage and wayfinding elements across all metro lines, in time for the renovation of the East Line and opening of the North-South Line. The new design is based on the already-existing R-net branding, though somewhat modified. It uses the Profile typeface and harks back to the original Unger design by using blue, white and red design elements. All of the wayfinding systems commissioned after the original 1974 one are designed by Mijksenaar.
Communications
A communications backbone for the Amsterdam Metro was installed by Thales, which is also responsible for associated maintenance.
Map
| Technology | Netherlands | null |
619020 | https://en.wikipedia.org/wiki/Mountain%20goat | Mountain goat | The mountain goat (Oreamnos americanus), also known as the Rocky Mountain goat, is a cloven-footed mammal that is endemic to the remote and rugged mountainous areas of western North America. A subalpine to truly alpine species, it is a sure-footed climber commonly seen on sheer rock faces, near-vertical cliffs and icy passages. Mountain goats generally avoid venturing down into lower elevations—except during seasonal food shortages or during particularly bad weather—as the extreme elevation which they inhabit is their primary defense against predators such as black and brown bears, pumas and wolves.
Despite its vernacular name and both genera being in the same subfamily (Caprinae), the mountain goat is not a member of Capra, the genus that includes all true goats (such as the wild goat (Capra aegagrus), from which the domestic goat is derived); rather, it is more closely allied with the other bovids known as “goat-antelopes”, including the European chamois (Rupicapra), the gorals (Naemorhedus), the takins (Budorcas) and the serows (Capricornis), of Japan and eastern South Asia.
Classification and evolution
The mountain goat is an even-toed ungulate of the order Artiodactyla and the family Bovidae (along with antelopes, gazelles, and cattle). It belongs to the subfamily Caprinae, along with true goats, wild sheep, the chamois, the muskox and other species. The takins of the Himalayan region, while not a sister lineage of the mountain goat, are nonetheless very closely related and almost coeval to the mountain goat; they evolved in parallel from an ancestral goat. Other members of this group are the bharal, the true goats, and the Himalayan tahr. The sheep lineage is also very closely related, while the muskox lineage is somewhat more distant. The mountain goats probably diverged from their relatives in the late Tortonian, some 7.5 to 8 million years ago.
Given that all major caprine lineages emerged in the Late Miocene and contain at least one but usually several species from the eastern Himalayan region, their most likely place of origin is between today's Tibet and Mongolia or nearby. The mountain goat's ancestors thus probably crossed the Bering Strait after they split from their relatives, presumably before the Wisconsinian glaciation. No Pliocene mountain goats have been identified yet; the known fossil record is fairly recent, entirely from North America, and barely differs from the living animals. In the Pleistocene era, the small prehistoric mountain goat Oreamnos harringtoni lived in the southern Rocky Mountains. Ancient DNA studies suggest that this was the sister species of the living mountain goat, not its ancestor; consequently, the living species would also date back to the Pleistocene at least.
The mountain goat is the only living species in the genus Oreamnos. The name Oreamnos is derived from the Greek term óros (stem ore-) meaning "mountain" (or, alternatively, oreas "mountain nymph") and the word amnós meaning "lamb".
General appearance and characteristics
Both male and female mountain goats have beards, short tails, and long black horns, in length, which contain yearly growth rings. They are protected from the elements by their woolly greyish white double coats. The fine, dense wool of their undercoats is covered by an outer layer of longer, hollow hairs. Mountain goats molt in spring by rubbing against rocks and trees, with the adult billies shedding their extra wool first and the pregnant nannies shedding last. Their coats help them to withstand winter temperatures as low as and winds of up to .
A male goat stands about at the shoulder to the waist and can weigh considerably more than the female (around 30% more in some cases). Male goats also have longer horns and longer beards than females. The head-and-body length can range from , with a small tail adding .
The mountain goat's feet are well-suited for climbing steep, rocky slopes with pitches exceeding 60°, with inner pads that provide traction and cloven hooves that can spread apart. The tips of their feet have sharp dewclaws that keep them from slipping. They have powerful shoulder and neck muscles that help propel them up steep slopes. Based on a field recording in the Rocky Mountains of Canada of a mountain goat climbing a 45-degree slope, researchers were able to measure the goat's whole body movement as it climbed. Researchers observed that when the goat propelled itself forward, it extended its back legs and the front legs were tucked close up to its chest during its first phase. During the second phase, the goat raised its back legs near to its chest, while the front leg's humerus stayed locked in a persistent location relative to the goat's chest, therefore allowing the elbow to be detained in close proximity to the whole body's center of balance. Extension of the elbow and carpal joints resulted in a vertical translation of the center of mass up the mountain slope.
Range and habitat
The mountain goat inhabits the Rocky Mountains and Cascade Range and other mountain regions of the Western Cordillera of North America, from Washington, Idaho and Montana through British Columbia and Alberta, into the southern Yukon and southeastern Alaska. British Columbia contains half of the world's population of mountain goats. Its northernmost range is said to be along the northern fringe of the Chugach Mountains in south-central Alaska. Introduced populations can also be found in such areas as Idaho, Wyoming, Utah, Nevada, Oregon, Colorado, South Dakota, and the Olympic Peninsula of Washington.
Mountain goats are the largest mammals found in their high-altitude habitats, which can exceed elevations of . They sometimes descend to sea level in coastal areas although they are primarily an alpine and subalpine species. The animals usually stay above the tree line throughout the year but they will migrate seasonally to higher or lower elevations within that range. Winter migrations to low-elevation mineral licks often take them several kilometers through forested areas.
Movement patterns
Daily movements by individual mountain goats are primarily confined to areas on the same mountain face, drainage basin, or alpine opening. Daily movements reflect an individual's needs for foraging, resting, thermoregulation and security from predators or disturbance. Seasonal movements primarily reflect nutritional needs (such as movements to and from mineral licks/salt lick), reproductive needs (in other words, movement of pre-parturient females to "kidding" areas; movement to rutting areas), and climatic influences (including movement to areas in response to foraging conditions). In general, seasonal movements are likely to exhibit a strong elevational component, whereby lower, forested elevations are used during the spring-summer (security cover effects) to access lower elevation mineral licks, and during winter (thermal cover effects) to access forage. The farthest movements are expected to be by dispersing mountain goats. Such movements are likely to involve mountain goats crossing forested valleys as they move between mountain blocks.
Diet
Mountain goats are herbivores and spend most of their time grazing. Their diets include grasses, herbs, sedges, ferns, mosses, lichens, and twigs and leaves from the low-growing shrubs and conifers of their high-altitude habitat.
In captivity, the mountain goat's diet can also include grain, alfalfa, fruits, vegetables and grass.
Life cycle and mating
In the wild, mountain goats usually live 12to15 years, with their lifespans limited by the wearing down of their teeth. In zoos, however, they can live for 16-20years.
Mountain goats reach sexual maturity at about 30 months. Nannies in a herd undergo synchronized estrus in late October through early December, at which time females and males participate in a mating ritual. Mature billies stare at nannies for long periods, dig rutting pits, and fight each other in showy (though occasionally dangerous) scuffles. Nannies often ignore young billies, who try to participate but are discounted in favor of older partners. Both females and males usually mate with multiple individuals during breeding season, although some billies try to keep other males away from certain nannies. After the breeding season is over, females and males move away from each other. Nannies form loose-knit nursery groups of up to 50 animals. The adult billies leave, often alone or with two or three other billies.
Kids are born in the spring (late May or early June) after a six-month gestation period. Nannies give birth, usually to a single offspring, after moving to an isolated ledge; post partum, they lick the kid dry and ingest the placenta. Kids weigh a little over at birth and begin to run and climb (or attempt to do so) within hours. Although lactation is mostly finished at one month, kids follow their mothers closely for the first year of life (or until the nanny gives birth again, if this does not occur the next breeding season); nannies protect their young by leading them out of danger, standing over them when faced by predators, and positioning themselves below their kids on steep slopes to stop freefalls.
Aggressive behavior
Nannies can be very competitive and protective of their space and food sources. They fight with one another for dominance in conflicts that can ultimately include all the nannies in the herd. In these battles, nannies circle each other with their heads lowered, displaying their horns. These conflicts can occasionally lead to injury or death, but are usually harmless. To avoid fighting, an animal may show a posture of nonaggression by stretching low to the ground.
In regions below the tree line, nannies use their fighting abilities to protect themselves and their offspring from predators. Predators, including wolves, wolverines, lynxes, and bears, attack goats of most ages given the opportunity. The cougar, or mountain lion, is perhaps the primary predator, being powerful enough to overwhelm the largest adults and uniquely nimble enough to navigate the rocky ecosystem of the goats. Though their size protects them from most potential predators in higher altitudes, nannies must sometimes defend their young from both bald and golden eagles, which can be a predatory threat to kids. Nannies have even been observed trying to dominate the more passive, but often heavier bighorn sheep that share some of their territory. In 2021, a mountain goat gored a grizzly bear to death in Yoho National Park, British Columbia.
Mountain goats introduced in the 1920s into Washington’s Olympic Mountains were in time found to be a nuisance there, in particular while seeking human urine and sweat for its salt content, the park lacking natural salt licks, and even aggressively approaching human visitors. One such goat killed a hiker in 2010. Officials finally chose to eradicate them from the Olympic Peninsula, removing hundreds, mostly by capturing them and relocating them to the Cascade Mountains.
Wool
Although mountain goats have never been domesticated and commercialized for their wool, pre-Columbian indigenous peoples of the Pacific Northwest Coast did incorporate their wool into their weaving by collecting spring moulted wool left by wild goats.
Threats
Avalanches can constitute 23-65% of Alaskan mountain goat mortality depending on the region.
| Biology and health sciences | Bovidae | Animals |
619201 | https://en.wikipedia.org/wiki/Nebulizer | Nebulizer | In medicine, a nebulizer (American English) or nebuliser (British English) is a drug delivery device used to administer medication in the form of a mist inhaled into the lungs. Nebulizers are commonly used for the treatment of asthma, cystic fibrosis, COPD and other respiratory diseases or disorders. They use oxygen, compressed air or ultrasonic power to break up solutions and suspensions into small aerosol droplets that are inhaled from the mouthpiece of the device. An aerosol is a mixture of gas and solid or liquid particles.
Medical uses
Guidelines
Various asthma guidelines, such as the Global Initiative for Asthma Guidelines [GINA], the British Guidelines on the management of Asthma, The Canadian Pediatric Asthma Consensus Guidelines, and United States Guidelines for Diagnosis and Treatment of Asthma each recommend metered dose inhalers in place of nebulizer-delivered therapies.
The European Respiratory Society acknowledge that although nebulizers are used in hospitals and at home they suggest much of this use may not be evidence-based.
Effectiveness
Recent evidence shows that nebulizers are no more effective than metered-dose inhalers (MDIs) with spacers. An MDI with a spacer may offer advantages to children who have acute asthma. Those findings refer specifically to the treatment of asthma and not to the efficacy of nebulisers generally, as for COPD for example. For COPD, especially when assessing exacerbations or lung attacks, there is no evidence to indicate that MDI (with a spacer) delivered medicine is more effective than administration of the same medicine with a nebulizer.
The European Respiratory Society highlighted a risk relating to droplet size reproducibility caused by selling nebulizer devices separately from nebulized solution. They found this practice could vary droplet size 10-fold or more by changing from an inefficient nebulizer system to a highly efficient one.
Two advantages attributed to nebulizers, compared to MDIs with spacers (inhalers), are their ability to deliver larger dosages at a faster rate, especially in acute asthma; however, recent data suggests actual lung deposition rates are the same. In addition, another trial found that a MDI (with spacer) had a lower required dose for clinical result compared to a nebulizer.
Beyond use in chronic lung disease, nebulizers may also be used to treat acute issues like the inhalation of toxic substances. One such example is the treatment of inhalation of toxic hydrofluoric acid (HF) vapors. Calcium gluconate is a first-line treatment for HF exposure to the skin. By using a nebulizer, calcium gluconate is delivered to the lungs as an aerosol to counteract the toxicity of inhaled HF vapors.
Aerosol deposition
The lung deposition characteristics and efficacy of an aerosol depend largely on the particle or droplet size. Generally, the smaller the particle the greater its chance of peripheral penetration and retention. However, for very fine particles below 0.5 μm in diameter there is a chance of avoiding deposition altogether and being exhaled. In 1966 the Task Group on Lung Dynamics, concerned mainly with the hazards of inhalation of environmental toxins, proposed a model for deposition of particles in the lung. This suggested that particles of more than 10 μm in diameter are most likely to deposit in the mouth and throat, for those of 5–10 μm diameter a transition from mouth to airway deposition occurs, and particles smaller than 5 μm in diameter deposit more frequently in the lower airways and are appropriate for pharmaceutical aerosols.
Nebulizing processes have been modeled using computational fluid dynamics.
Types
Pneumatic
Jet nebulizer
The most commonly used nebulizers are jet nebulizers, which are also called "atomizers". Jet nebulizers are connected by tubing to a supply of compressed gas, usually compressed air or oxygen to flow at high velocity through a liquid medicine to turn it into an aerosol that is inhaled by the patient. Currently there seems to be a tendency among physicians to prefer prescription of a pressurized Metered Dose Inhaler (pMDI) for their patients, instead of a jet nebulizer that generates a lot more noise (often 60 dB during use) and is less portable due to a greater weight. However, jet nebulizers are commonly used in hospitals for patients who have difficulty using inhalers, such as in serious cases of respiratory disease, or severe asthma attacks. The main advantage of the jet nebulizer is related to its low operational cost. If the patient needs to inhale medicine on a daily basis the use of a pMDI can be rather expensive. Today several manufacturers have also managed to lower the weight of the jet nebulizer to just over half a kilogram (just under one and a half pounds), and therefore started to label it as a portable device. Compared to all the competing inhalers and nebulizers, the noise and heavy weight is still the biggest draw back of the jet nebulizer.
Mechanical
Soft mist inhaler
The medical company Boehringer Ingelheim also invented a device named Respimat Soft Mist Inhaler in 1997. This new technology provides a metered dose to the user, as the liquid bottom of the inhaler is rotated clockwise 180 degrees by hand, adding a buildup tension into a spring around the flexible liquid container. When the user activates the bottom of the inhaler, the energy from the spring is released and imposes pressure on the flexible liquid container, causing liquid to spray out of 2 nozzles, thus forming a soft mist to be inhaled. The device features no gas propellant and no need for battery/power to operate. The average droplet size in the mist was measured to 5.8 micrometers, which could indicate some potential efficiency problems for the inhaled medicine to reach the lungs. Subsequent trials have proven this was not the case. Due to the very low velocity of the mist, the Soft Mist Inhaler in fact has a higher efficiency compared to a conventional pMDI. In 2000, arguments were launched towards the European Respiratory Society (ERS) to clarify/expand their definition of a nebulizer, as the new Soft Mist Inhaler in technical terms both could be classified as a "hand driven nebulizer" and a "hand driven pMDI".
Electrical
Ultrasonic wave nebulizer
Ultrasonic wave nebulizers were invented in 1965 as a new type of portable nebulizer. The technology inside an ultrasonic wave nebulizer is to have an electronic oscillator generate a high frequency ultrasonic wave, which causes the mechanical vibration of a piezoelectric element. This vibrating element is in contact with a liquid reservoir and its high frequency vibration is sufficient to produce a vapor mist via ultrasonic atomization. As they create aerosols from ultrasonic vibration instead of using a heavy air compressor, they only have a weight around . Another advantage is that the ultrasonic vibration is almost silent. Examples of these more modern type of nebulizers are: Omron NE-U17 and Beurer Nebulizer IH30.
Vibrating mesh technology
A new significant innovation was made in the nebulizer market around 2005, with creation of the ultrasonic Vibrating Mesh Technology (VMT). With this technology a mesh/membrane with 1000–7000 laser drilled holes vibrates at the top of the liquid reservoir, and thereby pressures out a mist of very fine droplets through the holes. This technology is more efficient than having a vibrating piezoelectric element at the bottom of the liquid reservoir, and thereby shorter treatment times are also achieved. The old problems found with the ultrasonic wave nebulizer, having too much liquid waste and undesired heating of the medical liquid, have also been solved by the new vibrating mesh nebulizers. Available VMT nebulizers include: Pari eFlow, Respironics i-Neb, Beurer Nebulizer IH50, and Aerogen Aeroneb. As the price of the ultrasonic VMT nebulizers is higher than models using previous technologies, most manufacturers continue to also sell the classic jet nebulizers.
Use and attachments
Nebulizers accept their medicine in the form of a liquid solution, which is often loaded into the device upon use. Corticosteroids and bronchodilators such as salbutamol (albuterol USAN) are often used, and sometimes in combination with ipratropium. The reason these pharmaceuticals are inhaled instead of ingested is in order to target their effect to the respiratory tract, which speeds onset of action of the medicine and reduces side effects, compared to other alternative intake routes.
Usually, the aerosolized medicine is inhaled through a tube-like mouthpiece, similar to that of an inhaler. The mouthpiece, however, is sometimes replaced with a face mask, similar to that used for inhaled anesthesia, for ease of use with young children or the elderly. Pediatric masks are often shaped like animals such as fish, dogs or dragons to make children less resistant to nebulizer treatments. Many nebulizer manufacturers also offer pacifier attachments for infants and toddlers. But mouthpieces are preferable if patients are able to use them since face-masks result in reduced lung delivery because of aerosol losses in the nose.
After use with corticosteroid, it is theoretically possible for patients to develop a yeast infection in the mouth (thrush) or hoarseness of voice (dysphonia), although these conditions are clinically very rare. To avoid these adverse effects, some clinicians suggest that the person who used the nebulizer should rinse his or her mouth. This is not true for bronchodilators; however, patients may still wish to rinse their mouths due to the unpleasant taste of some bronchodilating drugs.
History
The first "powered" or pressurized inhaler was invented in France by Sales-Girons in 1858. This device used pressure to atomize the liquid medication. The pump handle is operated like a bicycle pump. When the pump is pulled up, it draws liquid from the reservoir, and upon the force of the user's hand, the liquid is pressurized through an atomizer, to be sprayed out for inhalation near the user's mouth.
In 1864, the first steam-driven nebulizer was invented in Germany. This inhaler, known as "Siegle's steam spray inhaler", used the Venturi principle to atomize liquid medication, and this was the very beginning of nebulizer therapy. The importance of droplet size was not yet understood, so the efficacy of this first device was unfortunately mediocre for many of the medical compounds. The Siegle steam spray inhaler consisted of a spirit burner, which boiled water in the reservoir into steam that could then flow across the top and into a tube suspended in the pharmaceutical solution. The passage of steam drew the medicine into the vapor, and the patient inhaled this vapor through a mouthpiece made of glass.
The first pneumatic nebulizer fed from an electrically driven gas (air) compressor was invented in the 1930s and called a Pneumostat. With this device, a medical liquid (typically epinephrine chloride, used as a bronchial muscle relaxant to reverse constriction). As an alternative to the expensive electrical nebulizer, many people in the 1930s continued to use the much more simple and cheap hand-driven nebulizer, known as the Parke-Davis Glaseptic.
In 1956, a technology competing against the nebulizer was launched by Riker Laboratories (3M), in the form of pressurized metered-dose inhalers, with Medihaler-iso (isoprenaline) and Medihaler-epi (epinephrine) as the two first products. In these devices, the drug is cold-fill and delivered in exact doses through some special metering valves, driven by a gas propellant technology (i.e. Freon or a less environmentally damaging HFA).
In 1964, a new type of electronic nebulizer was introduced: the "ultrasonic wave nebulizer". Today the nebulizing technology is not only used for medical purposes. Ultrasonic wave nebulizers are also used in humidifiers, to spray out water aerosols to moisten dry air in buildings.
Some of the first models of electronic cigarettes featured an ultrasonic wave nebulizer (having a piezoelectric element vibrating and creating high-frequency ultrasound waves, to cause vibration and atomization of liquid nicotine) in combination with a vapouriser (built as a spray nozzle with an electric heating element). The most common type of electronic cigarettes currently sold, however, omit the ultrasonic wave nebulizer, as it was not found to be efficient enough for this kind of device. Instead, the electronic cigarettes now use an electric vaporizer, either in direct contact with the absorbent material in the "impregnated atomizer," or in combination with the nebulization technology related to a "spraying jet atomizer" (in the form of liquid droplets being out-sprayed by a high-speed air stream, that passes through some small venturi injection channels, drilled in a material absorbed with nicotine liquid).
| Biology and health sciences | General concepts_2 | Health |
619355 | https://en.wikipedia.org/wiki/Retainer%20%28orthodontics%29 | Retainer (orthodontics) | Orthodontic retainers are custom-made devices, usually made of wires or clear plastic, that hold teeth in position after surgery or any method of realigning teeth. Once a phase of orthodontic treatment has been completed to straighten teeth, there remains a lifelong risk of relapse (a tendency for teeth to return to their original position) due to a number of factors: recoil of periodontal fibres, pressure from surrounding soft tissues, the occlusion and patient’s continued growth and development. By using retainers to hold the teeth in their new position for a length of time, the surrounding periodontal fibres adapt to changes in the bone which can help minimize any changes to the final tooth position after the completion of orthodontic treatment. Retainers may also be used to treat overjets.
Removable retainers
Removable retainers include Hawley, Vacuum-formed, Begg and Barrer. They provide orthodontic retention when worn and they can be taken in and out of the mouth. They can be worn part-time or full-time if required or as advised by the orthodontist. In comparison to fixed retainers, removable retainers are easier to clean.
Hawley retainers
The best-known removable retainers are the Hawley retainers, which consist of a metal wires that typically surround the six anterior teeth and keep them in place. Hawley retainers are one of the oldest types of removable retainers. Named for its inventor, Dr. Charles A. Hawley, the labial wire, or Hawley bow, incorporates 2 omega loops for adjustment. It is anchored in an acrylic baseplate that sits in the palate (roof of the mouth). They are made from metal wire running along the outside of the teeth. There are many adaptations possible with Hawley retainers. The advantage of this type of retainer is that the metal wires can be adjusted to finish treatment and continue minor movement of the anterior teeth as needed. It also benefits from being robust and rigid, easy to construct and allows prosthetic tooth/teeth to be added onto with metal stops placed mesial and distal to the prosthetic teeth to prevent any relapse. To help fix rotations; acrylic facing can be added to the labial bow and a bite plane added to maintain the result of deep overbite correction. Also, to control the position of the canine, reverse U-loop can be employed. Additionally, to avoid occlusal interferences, the labial bow can be soldered to the cribs.
The main disadvantages of this type of retainer is its inferior aesthetics, interference with speech, risk of fracture and inferior retention of lower incisors in comparison to vacuum-formed retainers.
Recently, a more aesthetic version of the Hawley retainer has been developed. For this alternative, the front metal wire is replaced with a clear wire called the ASTICS. This retainer is intended to be adjustable similarly to the traditional Hawley retainer, which is not practical with vacuum-formed retainers.
Research shows that Hawley retainers are not effective for preventing incisor irregularity relapse. Hawley retainers also affect speech, especially the d, s, t, and i sounds, however as they are often only worn at night time, this concern may not be so prevalent. Research shows that participants that wear Hawley retainers report being more embarrassed about the appliance than wearing vacuum formed retainers and they found Hawley retainers more difficult to wear. However, if the patient has concerns with regards to the visible metal wire, a clear polyethylene bow can be used to enhance aesthetics.
The process of making an acrylic retainer was featured on the How It's Made TV show on the Science Channel on Season 23, Episode 16.
Vacuum-formed retainers
Another common type of removable retainer is the vacuum formed retainer (VFR). This is a polypropylene or polyvinylchloride (PVC) material. VFRs are made using a thermoforming process, using vacuum- or pressure-thermoforming. This clear or transparent retainer normally fits over the entire arch of teeth, however some designs fit only from canine to canine (clip-on retainer). The retainer is clear and so virtually invisible when worn. Hence, it can offer an aesthetic advantage relative to other retainers. VFRs, if worn 24 hours per day, do not allow the upper and lower teeth to touch, as the retainers cover the occlusal (biting) surfaces of the teeth. Some orthodontists feel that it is important for the top and bottom chewing surfaces to meet to allow for "favourable settling" to occur, thus the need for wearing the retainer only intermittently. VFRs are most commonly worn overnight and removed while eating. If worn while eating, they can behave as a reservoir enclosing the teeth with cariogenic substances and lead to decalcification of teeth over time (formation of cavities). The same can result if the retainer is inserted straight after a meal or drink. VFRs tend to be favoured as an orthodontic retainer as they are more aesthetic, interfere less with speech and are more economical than other removable retainers. However, for patients with disorders such as bruxism, VFRs can be prone to rapid breakage and deterioration due to grinding of the VFR against the opposing teeth. There is, however, an increased ease of fabrication over other retainers if they do break. Another advantage of VFRs is that there is evidence to suggest superior retention of lower incisors compared to other removable retainer types, but the best retention was achieved with fixed retainers.
Begg retainers
The Begg retainers have the labial bow extending around the distal aspect of the terminal molar. They allow occlusal settling, as no wire work crosses the occlusion. Begg retainers are robust enough to be worn during eating, however they are more retentive than Hawley retainers and the labial bow is less prone to distortion.
Barrer retainers
The Barrer retainers (aka Spring retainers) carry acrylic bows both labially and lingually. The original appliance extended only to the canines, however due to the risk of swallowing or aspiration, a modification which includes cribs on the first molars has been described. These retainers can be used to realign minor lower incisor relapse.
SMART retainers
Dental practitioners and orthodontists are sometimes resistant to providing younger patients with removable retainers due to the potential issues of the patient not wearing their retainer consistently after completing their orthodontic treatment. However, the recent innovation of SMART Retainers and related mobile applications that track, remind, and reward the patients for their continued use of their retainer may lead to a significant decrease in the number of adolescents that obtain additional orthodontic treatment in the future due to them not complying with their retainer usage. The mobile application has a unique feature that allows for parental notifications. If the child does not open their Smart Retainer Case, the parent can receive mobile notifications to help reduce retainer non-compliance.
Fixed retainers
Fixed retainers are often used to provide orthodontic retention and avoid relapse. They commonly consist of a wire bonded with acid etch and composite to the lingual/palatal surface of the anterior teeth. In fixed retainers, composite is usually placed to bond and to cover the wire, whilst ensuring no interference in the interdental space. Fixed retainers are used in situations where instability is more likely, such as severe rotations, periodontal disease and median diastemas. Occasionally the patient will require a removable retainer as well. Fixed bonded retainers can be designed with a smooth wire or flexible spiral wire which is also known as multi-strand wire. The most commonly used are multi-strand wire bonded to all six anterior teeth or a round stainless steel wire bonded to the canines only. Although fixed retainers depend less on the patient’s cooperation for regular wearing, they are more difficult to clean and therefore need more attention from the patient to prevent plaque accumulation and subsequent gingival inflammation.
Maintenance
Relapse can occur if retainers are not worn regularly post-treatment. The orthodontist may advise wearing the retainer for a set period of time or indefinitely after orthodontic treatment. Recent innovations of a Smart Retainer Case that utilizes: IoT, motion sensors, and a mobile application to track, remind, and reward patients, might be the most effective way of keeping post-orthodontic patients engaged with their final retention and continued wear of their removable retainer. If not cleaned properly, retainers can act as a food reservoir and lead to caries and gingival inflammation. Cleaning options for removable retainers include retainer/denture tablets, the use of a toothbrush and fragrance-free soap or non-abrasive toothpaste. If the water used to clean the retainer is too hot, it may cause shrinkage.
If the retainer is removable, it should be cleaned before inserting. If the retainer is fixed and cannot be removed, it should be cleaned during the nighttime home care routine. To clean a retainer use a wet toothbrush and gently scrub all surfaces of the retainer to remove any plaque and bacteria. A non-abrasive toothpaste should be used when cleaning a retainer. Alternatively, a mild soap or a brush dipped in mouthwash can also be used to clean it. When not in use keep the retainer in a dry container. Avoid leaving it in high temperatures or in the sun. Do not soak a retainer in liquids overnight. This will ensure the material will stay intact.
| Technology | Equipment | null |
619442 | https://en.wikipedia.org/wiki/Venule | Venule | A venule is a very small vein in the microcirculation that allows blood to return from the capillary beds to drain into the venous system via increasingly larger veins. Post-capillary venules are the smallest of the veins with a diameter of between 10 and 30 micrometres (μm). When the post-capillary venules increase in diameter to 50μm they can incorporate smooth muscle and are known as muscular venules. Veins contain approximately 70% of total blood volume, while about 25% is contained in the venules. Many venules unite to form a vein.
Structure
Post-capillary venules have a single layer of endothelium surrounded by a basal lamina. Their size is between 10 and 30 micrometers and are too small to contain smooth muscle. They are instead supported by pericytes that wrap around them. When the post-capillary venules increase in diameter to 50μm they can incorporate smooth muscle and are known as muscular venules. They have an inner endothelium composed of squamous endothelial cells that act as a membrane, a middle layer of muscle and elastic tissue and an outer layer of fibrous connective tissue. The middle layer is poorly developed so that venules have thinner walls than arterioles. They are porous so that fluid and blood cells can move easily from the bloodstream through their walls.
Short portal venules between the posterior pituitary and the anterior pituitary lobes provide an avenue for rapid hormonal exchange via the blood. Specifically within and between the pituitary lobes is anatomical evidence for confluent interlobe venules providing blood from the anterior to the neural lobe that would facilitate moment-to-moment sharing of information between lobes of the pituitary gland.
In contrast to regular venules, high endothelial venules are a special type of venule where the endothelium is made up of simple cuboidal cells. Lymphocytes exit the blood stream and enter the lymph nodes via these specialized venules when an infection is detected. Compared with arterioles, the venules are larger with much weaker muscular coat. They are the smallest united common branch in the human body.
| Biology and health sciences | Circulatory system | Biology |
619451 | https://en.wikipedia.org/wiki/Internal%20bleeding | Internal bleeding | Internal bleeding (also called internal haemorrhage) is a loss of blood from a blood vessel that collects inside the body, and is not usually visible from the outside. It can be a serious medical emergency but the extent of severity depends on bleeding rate and location of the bleeding (e.g. head, torso, extremities). Severe internal bleeding into the chest, abdomen, pelvis, or thighs can cause hemorrhagic shock or death if proper medical treatment is not received quickly. Internal bleeding is a medical emergency and should be treated immediately by medical professionals.
Signs and symptoms
Signs and symptoms of internal bleeding may vary based on location, presence of injury or trauma, and severity of bleeding. Common symptoms of blood loss may include:
Lightheadedness
Fatigue
Urinating less than usual
Confusion
Fast heart rate
Pale and/or cold skin
Thirst
Generalized weakness
Visible signs of internal bleeding include:
Blood in the urine
Dark black stools
Bright red stools
Bloody noses
Bruising
Throwing up blood
Of note, it is possible to have internal bleeding without any of the above symptoms, and pain may or may not be present.
A patient may lose more than 30% of their blood volume before there are changes in their vital signs or level of consciousness. This is called hemorrhagic or hypovolemic shock, which is a type of shock that occurs when there is not enough blood to reach organs in the body.
Causes
Internal bleeding can be caused by a broad number of things. We can break these up into three large categories:
Trauma, or direct injury to blood vessels within the body cavity
Genetic and acquired conditions, along with various medications, that result in an increased bleeding risk
Other
Traumatic
The most common cause of death in trauma is bleeding. Death from trauma accounts for 1.5 million of the 1.9 million deaths per year due to bleeding.
There are two types of trauma: penetrating trauma and blunt trauma.
Penetrating trauma is the most common cause of vascular injury and can result in internal bleeding. It can occur after a ballistic injury or stab wound. If penetrating trauma occurs in blood vessels close to the heart, it can quickly lead to hemorrhagic or hypovolemic shock, exsanguination, and death.
Blunt trauma is another cause of vascular injury that can result in internal bleeding. It can occur after a high speed deceleration in an automobile accident.
Non-traumatic
A number of pathological conditions and diseases can lead to internal bleeding. These include:
Blood vessel rupture as a result of high blood pressure, aneurysms, peptic ulcers, or ectopic pregnancy.
Other diseases linked to internal bleeding include cancer, hematologic disease, Vitamin K deficiency, and rare viral hemorrhagic fevers, such as the Ebola, Dengue or Marburg viruses.
Other
Internal bleeding could be a result of complications following surgery or other medical procedures. Some medications may also increase a person's risk for bleeding, such as anticoagulant drugs or antiplatelet drugs in the treatment of coronary artery disease.
Diagnosis
Vital signs
Blood loss can be estimated based on heart rate, blood pressure, respiratory rate, and mental status. Blood is circulated throughout the body and all major organ systems through a closed loop system. When there is damage to the blood vessel or the blood is thinner than the physiologic consistency, blood can exit the vessel which disrupts this close-looped system. The autonomic nervous system (ANS) responds in two large ways as an attempt to compensate for the opening in the system. These two actions are easily monitored by checking the heart rate and blood pressure. Blood pressure will initially decrease due to the loss of blood. This is where the ANS comes in and attempts to compensate by contracting the muscles that surround these vessels. As a result, a person who is bleeding internally may initially have a normal blood pressure. When the blood pressure falls below the normal range, this is called hypotension. The heart will start to pump faster causing the heart rate to increase, as an attempt to get blood delivered to vital organ systems faster. When the heart beats faster than the healthy and normal range, this is called tachycardia. If the bleeding is not controlled or stopped, a patient will experience tachycardia and hypotension, which altogether is a state of shock, called hemorrhagic shock.
Advanced trauma life support (ATLS) by the American College of Surgeons separates hemorrhagic shock into four categories.
Assessing circulation occurs after assessing the patient's airway and breathing (ABC (medicine)). If internal bleeding is suspected, a patient's circulatory system is assessed through palpation of pulses and doppler ultrasonography.
Physical examination
It is important to examine the person for visible signs that may suggest the presence of internal bleeding and/or the source of the bleed. Some of these signs may include:
a wound
bruising [ecchymosis]
blood collection [hematoma]
abnormal skin sensation [paresthesia]
signs of compartment syndrome
Imaging
If internal bleeding is suspected a FAST exam may be performed to look for bleeding in the abdomen.
If the patient has stable vital signs, they may undergo diagnostic imaging such as a CT scan. If the patient has unstable vital signs, they may not undergo diagnostic imaging and instead may receive immediate medical or surgical treatment.
Treatment
Management of internal bleeding depends on the cause and severity of the bleed. Internal bleeding is a medical emergency and should be treated immediately by medical professionals.
Fluid replacement
If a patient has low blood pressure (hypotension), intravenous fluids can be used until they can receive a blood transfusion. In order to replace blood loss quickly and with large amounts of IV fluids or blood, patients may need a central venous catheter. Patients with severe bleeding need to receive large quantities of replacement blood via a blood transfusion. As soon as the clinician recognizes that the patient may have a severe, continuing hemorrhage requiring more than 4 units in 1 hour or 10 units in 6 hours, they should initiate a massive transfusion protocol. The massive transfusion protocol replaces red blood cells, plasma, and platelets in varying ratios based on the cause of the bleeding (traumatic vs. non-traumatic).
Stopping the bleeding
It is crucial to stop the internal bleeding immediately (achieve hemostasis) after identifying its cause. The longer it takes to achieve hemostasis in people with traumatic causes (e.g. pelvic fracture) and non-traumatic causes (e.g. gastrointestinal bleeding, ruptured abdominal aortic aneurysm), the higher the death rate is.
Unlike with external bleeding, most internal bleeding cannot be controlled by applying pressure to the site of injury. Internal bleeding in the thorax and abdominal cavity (including both the intraperitoneal and retroperitoneal space) cannot be controlled with direct pressure (compression). A patient with acute internal bleeding in the thorax after trauma should be diagnosed, resuscitated, and stabilized in the Emergency Department in less than 10 minutes before undergoing surgery to reduce the risk of death from internal bleeding. A patient with acute internal bleeding in the abdomen or pelvis after trauma may require use of a REBOA device to slow the bleeding. The REBOA has also been used for non-traumatic causes of internal bleeding, including bleeding during childbirth and gastrointestinal bleeding.
Internal bleeding from a bone fracture in the arms or legs may be partially controlled with direct pressure using a tourniquet. After tourniquet placement, the patient may need immediate surgery to find the bleeding blood vessel.
Internal bleeding where the torso meets the extremities ("junctional sites" such as the axilla or groin) cannot be controlled with a tourniquet; however there is an FDA approved device known as an Abdominal Aortic and Junctional Tourniquet (AAJT) designed for proximal aortic control, although very few studies examining its use have been published. For bleeding at junctional sites, a dressing with a blood clotting agent (hemostatic dressing) should be applied.
A campaign is to improve the care of the bleeding known as Stop The Bleed campaign is also taking place.
| Biology and health sciences | Types | Health |
619594 | https://en.wikipedia.org/wiki/Aciclovir | Aciclovir | Aciclovir, also known as acyclovir, is an antiviral medication. It is primarily used for the treatment of herpes simplex virus infections, chickenpox, and shingles. Other uses include, prevention of cytomegalovirus infections following transplant, and severe complications of Epstein–Barr virus infection. It can be taken by mouth, applied as a cream, or injected.
Common side effects include nausea and diarrhea. Potentially serious side effects include kidney problems and low platelets. Greater care is recommended in those with poor liver or kidney function. It is generally considered safe for use in pregnancy with no harm having been observed. It appears to be safe during breastfeeding. Aciclovir is a nucleoside analogue that mimics guanosine. It works by decreasing the production of the virus's DNA.
Aciclovir was patented in 1974, by Burroughs Wellcome, and approved for medical use in 1981. It is on the World Health Organization's List of Essential Medicines. It is available as a generic medication and is marketed under many brand names worldwide. In 2022, it was the 134th most commonly prescribed medication in the United States, with more than 4million prescriptions.
Medical use
Aciclovir acts by inhibiting viral DNA replication and is used for the treatment of herpes simplex virus (HSV) and varicella zoster virus infections, including:
Genital herpes simplex (treatment and prevention)
Neonatal herpes simplex
Herpes simplex labialis (cold sores)
Shingles
Acute chickenpox in immunocompromised patients
Herpes simplex encephalitis
Acute mucocutaneous HSV infections in immunocompromised patients
Herpes of the eye and herpes simplex blepharitis (a chronic (long-term) form of herpes eye infection)
Prevention of herpes viruses in immunocompromised people (such as people undergoing cancer chemotherapy)
It has not been found to be effective against Epstein–Barr virus and its associated infectious mononucleosis. Aciclovir risks causing resistance to antiviral agents, and in 1% to 10% of cases can cause unpleasant side effects.
Aciclovir taken by mouth does not appear to decrease the risk of pain after shingles. In those with herpes of the eye, aciclovir may be more effective and safer than idoxuridine. It is unclear if aciclovir eye drops are more effective than brivudine eye drops.
Intravenous aciclovir is effective in treating severe medical conditions caused by different species of the herpes virus family, including severe localized infections of the herpes virus, severe genital herpes, chickenpox and herpesviral encephalitis. It is also effective in systemic or traumatic herpes infections, eczema herpeticum, and herpesviral meningitis.
Reviews of research dating from the 1980s show there is some effect in reducing the number and duration of lesions if aciclovir is applied at an early stage of an outbreak. Research shows effectiveness of topical aciclovir in both the early and late stages of the outbreak as well as improving methodologically and in terms of statistical certainty from previous studies. Aciclovir trials show that this agent has no role in preventing HIV transmission, but it can help slow HIV disease progression in people not taking anti-retroviral therapy (ART). This finding emphasizes the importance of testing simple, inexpensive non-ART strategies, such as aciclovir and cotrimoxazole, in people with HIV.
Pregnancy
The CDC and others have declared that during severe recurrent or first episodes of genital herpes, aciclovir may be used. For severe HSV infections (especially disseminated HSV), IV aciclovir may also be used.
Studies in mice, rabbits, and rats (with doses more than 10 times the equivalent of that used in humans) given during organogenesis have failed to demonstrate birth defects. Studies in rats in which they were given the equivalent to 63 times the standard steady-state human concentrations of the drug on day 10 of gestation showed head and tail anomalies.
Aciclovir is recommended by the CDC for treatment of varicella during pregnancy, especially during the second and third trimesters.
Aciclovir is excreted in breast milk, therefore it is recommended that caution should be used in breast-feeding women. It has been shown in limited test studies that the nursing infant is exposed to approximately 0.3 mg/kg/day following oral administration of aciclovir to the mother. If nursing mothers have herpetic lesions near or on the breast, breast-feeding should be avoided.
Adverse effects
Systemic therapy
Common adverse drug reactions (≥1% of patients) associated with systemic aciclovir therapy (oral or IV) include nausea, vomiting, diarrhea, encephalopathy (with IV use only), injection site reactions (with IV use only) and headache. In high doses, hallucinations have been reported. Infrequent adverse effects (0.1–1% of patients) include agitation, vertigo, confusion, dizziness, oedema, arthralgia, sore throat, constipation, abdominal pain, hair loss, rash and weakness. Rare adverse effects (<0.1% of patients) include coma, seizures, neutropenia, leukopenia, crystalluria, anorexia, fatigue, hepatitis, Stevens–Johnson syndrome, toxic epidermal necrolysis, thrombotic thrombocytopenic purpura, anaphylaxis, and Cotard's syndrome.
Intravenous aciclovir may cause reversible nephrotoxicity in up to 5% to 10% of patients because of precipitation of aciclovir crystals in the kidney. Aciclovir crystalline nephropathy is more common when aciclovir is given as a rapid infusion and in patients with dehydration and preexisting renal impairment. Adequate hydration, a slower rate of infusion, and dosing based on renal function may reduce this risk.
The aciclovir metabolite 9-Carboxymethoxymethylguanine (9-CMMG) has been shown to play a role in neurological adverse events, particularly in older people and those with reduced renal function.
Topical therapy
Aciclovir topical cream is commonly associated (≥1% of patients) with dry or flaking skin or transient stinging/burning sensations. Infrequent adverse effects include erythema or itch. When applied to the eye, aciclovir is commonly associated (≥1% of patients) with transient mild stinging. Infrequently (0.1–1% of patients), ophthalmic aciclovir is associated with superficial punctate keratitis or allergic reactions.
Drug interactions
Ketoconazole: In-vitro replication studies have found a synergistic, dose-dependent antiviral activity against HSV-1 and HSV-2 when given with aciclovir. However, this effect has not been clinically established and more studies need to be done to evaluate the true potential of this synergy.
Probenecid: Reports of increased half-life of aciclovir, as well as decreased urinary excretion and renal clearance have been shown in studies where probenecid is given simultaneously with aciclovir.
Interferon: Synergistic effects when administered with aciclovir and caution should be taken when administering aciclovir to patients receiving IV interferon.
Zidovudine: Although administered often with aciclovir in HIV patients, neurotoxicity has been reported in at least one patient who presented with extreme drowsiness and lethargy 30–60 days after receiving IV aciclovir; symptoms resolved when aciclovir was discontinued.
Detection in biological fluids
Aciclovir may be quantitated in plasma or serum to monitor for drug accumulation in patients with renal dysfunction or to confirm a diagnosis of poisoning in acute overdose victims.
Mechanism of action
Aciclovir is converted by viral thymidine kinase to aciclovir monophosphate, which is then converted by host cell kinases to aciclovir triphosphate (ACV-TP, also known as aciclo-GTP). ACV-TP is a very potent inhibitor of viral DNA replication. ACV-TP competitively inhibits and inactivates the viral DNA polymerase. Its monophosphate form also incorporates into the viral DNA, resulting in chain termination.
Resistance
Resistance to aciclovir is rare in people with healthy immune systems but is more common (up to 10%) in people with immunodeficiencies on chronic antiviral prophylaxis (transplant recipients, people with acquired immunodeficiency syndrome due to HIV infection). Mechanisms of resistance in HSV include deficient viral thymidine kinase; and mutations to viral thymidine kinase or DNA polymerase, altering substrate sensitivity.
Microbiology
Aciclovir is active against most species in the herpesvirus family. In descending order of activity:
Herpes simplex virus type I (HSV-1)
Herpes simplex virus type II (HSV-2)
Varicella zoster virus
Epstein–Barr virus
Human cytomegalovirus – least activity
Pharmacokinetics
Aciclovir is poorly water-soluble and has poor oral bioavailability (15–30%), hence intravenous administration is necessary if high concentrations are required. When orally administered, peak plasma concentration occurs after 1–2 hours. According to the Biopharmaceutical Classification System, aciclovir is a Class III drug, i.e., soluble with low intestinal permeability. Aciclovir has a high distribution rate; protein binding is reported to range from 9 to 33%. The elimination half-life (t1/2) of aciclovir depends according to age group; neonates have a t1/2 of 4 hours, children 1–12 years have a t1/2 of 2–3 hours whereas adults have a t1/2 of 3 hours.
Chemistry
Details of the synthesis of aciclovir were first published by scientists from the University at Buffalo.
In the first step shown, 2,6-dichloropurine was alkylated with 1-benzoyloxy-2-chloromethoxyethane. The chlorine group at the 6-position of the heterocyclic ring is more reactive than the chlorine at the 2-position, hence it can be selectively replaced by an amino group, which was then converted to an amide using nitrous acid. Finally, the remaining chlorine was replaced by the amino group of aciclovir using ammonia in methanol. This synthesis and other methods for preparing the compound have been reviewed.
History
Aciclovir was seen as the start of a new era in antiviral therapy, as it is extremely selective and low in cytotoxicity. Since the discovery in the mid-1970s, it has been used as an effective drug for the treatment of infections caused by most known species of the herpesvirus family, including herpes simplex and varicella zoster viruses. Nucleosides isolated from a Caribbean sponge, Cryptotethya crypta, were the basis for the synthesis of aciclovir. It was codiscovered by Howard Schaeffer following his work with Robert Vince, S. Bittner and S. Gurwara on the adenosine analog acycloadenosine which showed promising antiviral activity. Later, Schaeffer joined Burroughs Wellcome and continued the development of aciclovir with pharmacologist Gertrude B. Elion. A U.S. patent on aciclovir listing Schaeffer as inventor was issued in 1979. Vince later invented abacavir, an nRTI drug for HIV patients. Elion was awarded the 1988 Nobel Prize in Medicine, partly for the development of aciclovir.
A related prodrug form, valaciclovir came into medical use in 1995. It is converted to aciclovir in the body after absorption.
In 2009, acyclovir in combination with hydrocortisone cream, marketed as Xerese, was approved in the United States for the early treatment of recurrent herpes labialis (cold sores) to reduce the likelihood of ulcerative cold sores and to shorten the lesion healing time in adults and children (six years of age and older).
Society and culture
Names
It was originally marketed as Zovirax; patents expired in the 1990s and since then it has been generic and is marketed under many brand names worldwide.
| Biology and health sciences | Antiviral drugs | Health |
619602 | https://en.wikipedia.org/wiki/Genitourinary%20system | Genitourinary system | The genitourinary system, or urogenital system, are the sex organs of the reproductive system and the organs of the urinary system. These are grouped together because of their proximity to each other, their common embryological origin and the use of common pathways. Because of this, the systems are sometimes imaged together. In placental mammals (including humans), the male urethra goes through and opens into the penis while the female urethra and vagina empty through the vulva.
The term "apparatus urogenitalis" was used in Nomina Anatomica (under splanchnologia) but is not used in the current Terminologia Anatomica.
Development
The urinary and reproductive organs are developed from the intermediate mesoderm. The permanent organs of the adult are preceded by a set of structures that are purely embryonic and that, with the exception of the ducts, disappear almost entirely before the end of fetal life. These embryonic structures are on either side: the pronephros, the mesonephros and the metanephros of the kidney, and the Wolffian and Müllerian ducts of the sex organ. The pronephros disappears very early; the structural elements of the mesonephros mostly degenerate, but the gonad is developed in their place, with which the Wolffian duct remains as the duct in males, and the Müllerian as that of the female. Some of the tubules of the mesonephros form part of the permanent kidney.
Structures
Urethra
Female Urethra
The urethra of an adult human female is 3-4 cm long. The female urethra is located between the bladder neck to the external urethral orifice and is behind the symphysis pubis.The urethral wall is composed of an inner epithelial lining, a sub-mucosa layer containing vascular supply, a thin fascial layer, and two layers of smooth muscle.
Male Urethra
The urethra of an adult human male is 18-20 cm long. It has a diameter of 8-9 mm.The male urethra is divided into two sections.
Disorders
Disorders of the genitourinary system includes a range of disorders from those that are asymptomatic to those that manifest an array of signs and symptoms. Causes for these disorders include congenital anomalies, infectious diseases, trauma, or conditions that secondarily involve the urinary structure.
To gain access to the body, pathogens can penetrate mucous membranes lining the genitourinary tract.
Malformations
Urogenital malformations include:
Hypospadias
Epispadias
Labial fusion
Varicocele
As a medical specialty, genitourinary pathology is the subspecialty of surgical pathology which deals with the diagnosis and characterization of neoplastic and non-neoplastic diseases of the urinary tract, male genital tract and testes. However, medical disorders of the kidneys are generally within the expertise of renal pathologists. Genitourinary pathologists generally work closely with urologic surgeons.
| Biology and health sciences | Reproductive system | Biology |
619798 | https://en.wikipedia.org/wiki/Club%20%28weapon%29 | Club (weapon) | A club (also known as a cudgel, baton, bludgeon, truncheon, cosh, nightstick, or impact weapon) is a short staff or stick, usually made of wood, wielded as a weapon or tool since prehistory. There are several examples of blunt-force trauma caused by clubs in the past, including at the site of Nataruk in Turkana, Kenya, described as the scene of a prehistoric conflict between bands of hunter-gatherers 10,000 years ago.
Most clubs are small enough to be swung with one hand, although larger clubs may require the use of two to be effective. Various specialized clubs are used in martial arts and other fields, including the law-enforcement baton. The military mace is a more sophisticated descendant of the club, typically made of metal and featuring a spiked, knobbed, or flanged head attached to a shaft.
Examples of cultural depictions of clubs may be found in mythology, where they are associated with strong figures such as Hercules or the Japanese oni, or in popular culture, where they are associated with primitive cultures, especially cavemen. Ceremonial maces may also be displayed as a symbol of governmental authority.
The wounds inflicted by a club are generally known as strike trauma or blunt-force trauma injuries.
Law enforcement
Police forces and their predecessors have traditionally favored the use, whenever possible, of less lethal weapons than guns or blades. Until recent times, when alternatives such as tasers and capsicum spray became available, this category of policing weapon has generally been filled by some form of wooden club variously termed a truncheon, baton, nightstick, or lathi. Short, flexible clubs are also often used, especially by plainclothes officers who need to avoid notice. These are known colloquially as blackjacks, saps, or coshes.
Conversely, criminals have been known to arm themselves with an array of homemade or improvised clubs, generally of easily concealable sizes, or which can be explained as being carried for legitimate purposes (such as baseball bats).
In addition, Shaolin monks and members of other religious orders around the world have employed cudgels from time to time as defensive weapons.
Types
Though perhaps the simplest of all weapons, clubs come in many varieties, including:
Aklys – a club with an integrated leather thong, used to return it to the hand after snapping it at an opponent. Used by the legions of the Roman Empire.
Ball club – These clubs were used by Native Americans. There are two types; the stone ball clubs that were used mostly by early Plains, Plateau and Southwest Native Indians and the wooden ball clubs that the Huron and Iroquois tribes used. These consisted of a relatively free-moving head of rounded stone or wood attached to a wooden handle.
Bang – Chinese military weapon type used in medieval times. Also used in modern Wushu showcase and martial-arts practice.
Baseball, cricket and T-ball bats – The baseball bat is often used as an improvised weapon, much like the pickaxe handle. In countries where baseball is not commonly played, baseball bats are often first thought of as weapons. Tee ball bats are also used in this manner. Their smaller size and lighter weight make the bats easier to handle in one hand than a baseball bat. Cricket bats are heavier and their flat shape and short handle make them unwieldy as weapons, but they are more commonly available than baseball bats in some countries.
Baton or truncheon – forms used by law enforcement.
Blackjack or cosh – a weighted club designed to stun the subject.
Bian – a tubular club used by medieval Chinese infantry and generals.
Clava (full name clava mere okewa) – a traditional stone hand-club used by Mapuche Indians in Chile, featuring a long flat body. In Spanish, it is known as clava cefalomorfa. It has some ritual importance as a special sign of distinction carried by the tribal chief.
Cudgel – A stout stick carried by peasants during the Middle Ages. It functioned as a walking staff and a weapon for both self-defence and wartime. Clubmen revolted in several localities against the excesses of soldiers on both sides during the English Civil War. During the 18th century singlestick fighting (a training sport for the use of the single handed backsword) was called singlesticking, or cudgel-play.
Crowbar – a tool commonly used as an improvised weapon, though some examples are too large to be wielded with a single hand, and therefore should be classified as staves.
Flashlight – A large metal flashlight, such as a Maglite, can make a very effective improvised club. Though not specifically classified as a weapon, it is often carried for self-defense by security guards, bouncers and civilians, especially in countries where carrying weapons is restricted.
Gata – a Fijian war club
Ghioagă – a Romanian club similar to a shillelagh; also called Bâtă (the name comes from Latin batt(u)ere – battery). This was used as a weapon in group fights against Ottoman Empire by irregular troops made up of peasants, vassals to local Princes in Wallachia and Moldavia. Early mentions of it occur from the 15th century in some historical sources.
Gunstock war club – a war club stylized as the butt of a rifle
Jiǎn – a type of quad-edged straight club specifically designed to break other weapons with sharp edges.
Jutte or jitte – a distinctive weapon of the samurai police, consisting of an iron rod with a hook. It could parry and disarm a sword-wielding assailant without serious injury. Eventually, the jutte also came to be considered a symbol of official status.
Kanabō (nyoibo, konsaibo, tetsubō, ararebo) – Various types of different-sized Japanese clubs made of wood and or iron, usually with iron spikes or studs. First used by the samurai.
Kanak war clubs – traditional maces used by the Kanak people of New Caledonia
Kiyoga – a spring baton similar in concept to the Asp collapsible police baton, but with the center section made of a heavy-duty steel spring. The tip and first section slide into the spring, and the whole nests into a seven-inch handle. To deploy the kiyoga, all that is necessary is to grasp the handle and swing. This causes the parts to extend from the handle into a baton seventeen inches long. The kiyoga has one advantage over a conventional collapsible baton: it can reach around a raised arm trying to block it to strike the head.
Knobkerrie – a war club of southern and eastern Africa with a distinctive knob on the end
Kubotan – a short, thin, lightweight club often used by law enforcement officers, generally to apply pressure against selected points of the body in order to encourage compliance without inflicting injury.
Leangle – an Australian Aboriginal fighting-club with a hooked striking head, typically nearly at right angles to the weapon's shaft. The name comes from Kulin languages such as Wemba-Wemba and Woiwurrung, based on the word lia (tooth).
Life preserver (also hyphenated life-preserver) – a short, often weighted club intended for self-defense. Mentioned in Gilbert and Sullivan's 1879 comic opera The Pirates of Penzance and in several Sherlock Holmes stories.
Mace – a metal club with a heavy head on the end, designed to deliver very powerful blows. The head of a mace may also have small studs forged into it. The mace is often confused with the spiked morning star or with the articulated flail.
Mere – short, broad-bladed Māori club, usually made from nephrite jade and used for making forward-striking thrusts
Morning star – a medieval club-like weapon consisting of a shaft with an attached ball adorned with one or more spikes
Nulla-nulla – a short, curved hardwood club, used as a hunting weapon and in tribal in-fighting, by the Aboriginal people of Australia
Nunchaku (also called nunchucks) – an Asian weapon consisting of two clubs, connected by a short rope, thong or chain, and usually used with one club in hand and the other swung as a flail.
Oslop – a two-handed, very heavy, often iron-shod, Russian club that was used as the cheapest and the most readily available infantry weapon.
Paddle club – common in the Solomon Islands, these clubs could be used in warfare or for propelling a small dugout canoe.
Pickaxe handle – the (usually wooden) haft of a pickaxe used as a club
Racket (sports equipment)
Rungu (Swahili, plural marungu) – a wooden throwing club or baton bearing special symbolism and significance in certain East African tribal cultures. It is especially associated with Maasai morans (male warriors) who have traditionally used it in warfare and for hunting.
Sali, a Fijian war club
Sally rod – a long, thin wooden stick, generally made from willow (Latin salix), and used chiefly in the past in Ireland as a disciplinary implement, but also sometimes used like a club (without the fencing-like technique of stick fighting) in fights and brawls. In Japan this type of stick is called the Hanbō meaning half stick, and in FMA (Filipino martial arts) it is called the eskrima or escrima stick, often made from rattan.
Shillelagh – a wooden club or cudgel, typically made from a stout knotty stick with a large knob on the end, that is associated with Ireland in folklore
Slapjack – a variation of the blackjack consisting of a longer strap which lets it be used like a flail, and can be used as a club or for trapping techniques as seen in the use of nunchaku and other flexible weapons
Supi – a war club of the Solomon Islands
Telescopic baton – a rigid baton capable of collapsing to a shorter length for greater portability and concealability
Tipstaff – a ceremonial rod used by a court officer of the same name
Tonfa or side-handle baton – a club of Okinawan origin featuring a second handle mounted perpendicular to the shaft
Totokia – a Fijian spiked club
Trench raiding club – a type of melee weapon used by both sides in World War I
Ula – traditional throwing club from Fiji
U'u – an exquisitely-carved ceremonial club from the Marquesan Islands, used as a chiefly status symbol
Waddy – a heavy hardwood club, used as a weapon for hunting and in tribal in-fighting, and also as a tool, by the Aboriginal people of Australia. The word waddy describes a club from New South Wales, but Australians also use the word generally to include other Aboriginal clubs, including the nulla nulla and leangle.
Worraga – An Australian-aboriginal club with boomerang-like aerodynamics. Can be thrown or hand-held.
Animal appendages
Some animals have limbs or appendages resembling clubs, such as:
Ankylosaurus (armored dinosaur)
Anodontosaurus (armored dinosaur)
Club-winged manakin (extant bird)
Dyoplosaurus (armored dinosaur)
Jamaican ibis (extinct bird)
Mantis shrimp (marine crustacean)
Nodocephalosaurus (armored dinosaur)
Rodrigues solitaire (extinct bird)
Talarurus (armored dinosaur)
Gallery
| Technology | Melee weapons | null |
619926 | https://en.wikipedia.org/wiki/Gravitational%20collapse | Gravitational collapse | Gravitational collapse is the contraction of an astronomical object due to the influence of its own gravity, which tends to draw matter inward toward the center of gravity. Gravitational collapse is a fundamental mechanism for structure formation in the universe. Over time an initial, relatively smooth distribution of matter, after sufficient accretion, may collapse to form pockets of higher density, such as stars or black holes.
Star formation involves a gradual gravitational collapse of interstellar medium into clumps of molecular clouds and potential protostars. The compression caused by the collapse raises the temperature until thermonuclear fusion occurs at the center of the star, at which point the collapse gradually comes to a halt as the outward thermal pressure balances the gravitational forces. The star then exists in a state of dynamic equilibrium. During the star's evolution a star might collapse again and reach several new states of equilibrium.
Star formation
An interstellar cloud of gas will remain in hydrostatic equilibrium as long as the kinetic energy of the gas pressure is in balance with the potential energy of the internal gravitational force. Mathematically this is expressed using the virial theorem, which states that to maintain equilibrium, the gravitational potential energy must equal twice the internal thermal energy. If a pocket of gas is massive enough that the gas pressure is insufficient to support it, the cloud will undergo gravitational collapse. The critical mass above which a cloud will undergo such collapse is called the Jeans mass. This mass depends on the temperature and density of the cloud but is typically thousands to tens of thousands of solar masses.
Stellar remnants
At what is called the star's death (when a star has burned out its fuel supply), it will undergo a contraction that can be halted only if it reaches a new state of equilibrium. Depending on the mass during its lifetime, these stellar remnants can take one of three forms:
White dwarfs, in which gravity is opposed by electron degeneracy pressure
Neutron stars, in which gravity is opposed by neutron degeneracy pressure and short-range repulsive neutron–neutron interactions mediated by the strong force
Black hole, in which there is no force strong enough to resist gravitational collapse
White dwarf
The collapse of the stellar core to a white dwarf takes place over tens of thousands of years, while the star blows off its outer envelope to form a planetary nebula. If it has a companion star, a white dwarf-sized object can accrete matter from the companion star. Before it reaches the Chandrasekhar limit (about one and a half times the mass of the Sun, at which point gravitational collapse would start again), the increasing density and temperature within a carbon-oxygen white dwarf initiate a new round of nuclear fusion, which is not regulated because the star's weight is supported by degeneracy rather than thermal pressure, allowing the temperature to rise exponentially. The resulting runaway carbon detonation completely blows the star apart in a type Ia supernova.
Neutron star
Neutron stars are formed by the gravitational collapse of the cores of larger stars. They are the remnant of supernova types Ib, Ic, and II. Neutron stars are expected to have a skin or "atmosphere" of normal matter on the order of a millimeter thick, underneath which they are composed almost entirely of closely packed neutrons called neutron matter with a slight dusting of free electrons and protons mixed in. This degenerate neutron matter has a density of about .
The appearance of stars composed of exotic matter and their internal layered structure is unclear since any proposed equation of state of degenerate matter is highly speculative. Other forms of hypothetical degenerate matter may be possible, and the resulting quark stars, strange stars (a type of quark star), and preon stars, if they exist, would, for the most part, be indistinguishable from a neutron star: In most cases, the exotic matter would be hidden under a crust of "ordinary" degenerate neutrons.
Black holes
According to Einstein's theory, for even larger stars, above the Landau–Oppenheimer–Volkoff limit, also known as the Tolman–Oppenheimer–Volkoff limit (roughly double the mass of the Sun) no known form of cold matter can provide the force needed to oppose gravity in a new dynamical equilibrium. Hence, the collapse continues with nothing to stop it.
Once a body collapses to within its Schwarzschild radius it forms what is called a black hole, meaning a spacetime region from which not even light can escape. It follows from general relativity and the theorem of Roger Penrose that the subsequent formation of some kind of singularity is inevitable. Nevertheless, according to Penrose's cosmic censorship hypothesis, the singularity will be confined within the event horizon bounding the black hole, so the spacetime region outside will still have a well-behaved geometry, with strong but finite curvature, that is expected to evolve towards a rather simple form describable by the historic Schwarzschild metric in the spherical limit and by the more recently discovered Kerr metric if angular momentum is present. If the precursor has a magnetic field, it is dispelled during the collapse, as black holes are thought to have no magnetic field of their own.
On the other hand, the nature of the kind of singularity to be expected inside a black hole remains rather controversial. According to theories based on quantum mechanics, at a later stage, the collapsing object will reach the maximum possible energy density for a certain volume of space or the Planck density (as there is nothing that can stop it). This is the point at which it has been hypothesized that the known laws of gravity cease to be valid. There are competing theories as to what occurs at this point. For example loop quantum gravity predicts that a Planck star would form. Regardless, it is argued that gravitational collapse ceases at that stage and a singularity, therefore, does not form.
Theoretical minimum radius for a star
The radii of larger mass neutron stars (about 2.8 solar mass) are estimated to be about 12 km, or approximately 2 times their equivalent Schwarzschild radius.
It might be thought that a sufficiently massive neutron star could exist within its Schwarzschild radius (1.0 SR) and appear like a black hole without having all the mass compressed to a singularity at the center; however, this is probably incorrect. Within the event horizon, the matter would have to move outward faster than the speed of light in order to remain stable and avoid collapsing to the center. No physical force, therefore, can prevent a star smaller than 1.0 SR from collapsing to a singularity (at least within the currently accepted framework of general relativity; this does not hold for the Einstein–Yang–Mills–Dirac system). A model for the nonspherical collapse in general relativity with the emission of matter and gravitational waves has been presented.
| Physical sciences | Celestial mechanics | Astronomy |
619980 | https://en.wikipedia.org/wiki/Thermocline | Thermocline | A thermocline (also known as the thermal layer or the metalimnion in lakes) is
a distinct layer based on temperature within a large body of fluid (e.g. water, as in an ocean or lake; or air, e.g. an atmosphere) with a high gradient of distinct temperature differences associated with depth. In the ocean, the thermocline divides the upper mixed layer from the calm deep water below.
Depending largely on season, latitude, and turbulent mixing by wind, thermoclines may be a semi-permanent feature of the body of water in which they occur, or they may form temporarily in response to phenomena such as the radiative heating/cooling of surface water during the day/night. Factors that affect the depth and thickness of a thermocline include seasonal weather variations, latitude, and local environmental conditions, such as tides and currents.
Oceans
Most of the heat energy of the sunlight that strikes the Earth is absorbed in the first few centimeters at the ocean's surface, which heats during the day and cools at night as heat energy is lost to space by radiation. Waves mix the water near the surface layer and distribute heat to deeper water such that the temperature may be relatively uniform in the upper , depending on wave strength and the existence of surface turbulence caused by currents. Below this mixed layer, the temperature remains relatively stable over day/night cycles. The temperature of the deep ocean drops gradually with depth. As saline water does not freeze until it reaches (colder as depth and pressure increase) the temperature well below the surface is usually not far from zero degrees.
The thermocline varies in depth. It is semi-permanent in the tropics, variable in temperate regions and shallow to nonexistent in the polar regions, where the water column is cold from the surface to the bottom. A layer of sea ice will act as an insulation blanket. The first accurate global measurements were made during the oceanographic expedition of HMS Challenger.
In the open ocean, the thermocline is characterized by a negative sound speed gradient, making the thermocline important in submarine warfare because it can reflect active sonar and other acoustic signals. This stems from a discontinuity in the acoustic impedance of water created by the sudden change in density.
In scuba diving, a thermocline where water drops in temperature by a few degrees Celsius quite suddenly can sometimes be observed between two bodies of water, for example where colder upwelling water runs into a surface layer of warmer water. It gives the water an appearance of wrinkled glass, the kind often used in bathroom windows to obscure the view, and is caused by the altered refractive index of the cold or warm water column. These same schlieren can be observed when hot air rises off the tarmac at airports or desert roads and is the cause of mirages.
Thermocline seasonality
The thermocline in the ocean can vary in depth and strength seasonally. This is particularly noticeable in mid-latitudes with a thicker mixed layer in the winter and thinner mixed layer in summer. The cooler winter temperatures cause the thermocline to drop to further depths and warmed summer temperatures bring the thermocline back to the upper layer. In areas around the tropics and subtropics, the thermocline may become even thinner in the summer than in other locations. At higher latitudes, around the poles, there is more of a seasonal thermocline than a permanent one with warmer surface waters. This is where there is a dichothermal layer instead.
In the Northern hemisphere, the maximum temperatures at the surface occur through August and September and minimum temperatures occur through February and March with total heat content being lowest in March. This is when the seasonal thermocline starts to build back up after being broken down through the colder months.
A permanent thermocline is one that is not affected by season and lies below the yearly mixed layer maximum depth.
Other water bodies
Thermoclines can also be observed in lakes. In colder climates, this leads to a phenomenon called stratification. During the summer, warm water, which is less dense, will sit on top of colder, denser, deeper water with a thermocline separating them. The warm layer is called the epilimnion and the cold layer is called the hypolimnion. Because the warm water is exposed to the sun during the day, a stable system exists and very little mixing of warm water and cold water occurs, particularly in calm weather.
One result of this stability is that as the summer wears on, there is less and less oxygen below the thermocline as the water below the thermocline never circulates to the surface and organisms in the water deplete the available oxygen. As winter approaches, the temperature of the surface water will drop as nighttime cooling dominates heat transfer. A point is reached where the density of the cooling surface water becomes greater than the density of the deep water and overturning begins as the dense surface water moves down under the influence of gravity. This process is aided by wind or any other process (currents for example) that agitates the water. This effect also occurs in Arctic and Antarctic waters, bringing water to the surface which, although low in oxygen, is higher in nutrients than the original surface water. This enriching of surface nutrients may produce blooms of phytoplankton, making these areas productive.
As the temperature continues to drop, the water on the surface may get cold enough to freeze and the lake/ocean begins to ice over. A new thermocline develops where the densest water () sinks to the bottom, and the less dense water (water that is approaching the freezing point) rises to the top. Once this new stratification establishes itself, it lasts until the water warms enough for the 'spring turnover,' which occurs after the ice melts and the surface water temperature rises to 4 °C. During this transition, a thermal bar may develop.
Waves can occur on the thermocline, causing the depth of the thermocline as measured at a single location to oscillate (usually as a form of seiche). Alternately, the waves may be induced by flow over a raised bottom, producing a thermocline wave which does not change with time, but varies in depth as one moves into or against the flow.
Atmosphere
The thermal boundary between the troposphere (lower atmosphere) and the stratosphere (upper atmosphere) is a thermocline. Temperature generally decreases with altitude, but the heat from the day's exposure to sun is released at night, which can create a warm region at ground with colder air above. This is known as an inversion (a further example of a thermocline). At sunrise, the sun's energy warms the ground, causing the warming air to rise, thus destabilizing and eventually reversing the inversion layer.
This phenomenon was first applied to the field of noise pollution study in the 1960s, contributing to the design of urban highways and noise barriers.
| Physical sciences | Earth science basics: General | Earth science |
619985 | https://en.wikipedia.org/wiki/Menai%20Bridge | Menai Bridge | Menai Bridge (; usually referred to colloquially as Y Borth) is a town and community on the Isle of Anglesey in north-west Wales. It overlooks the Menai Strait and lies by the Menai Suspension Bridge, built in 1826 by Thomas Telford, just over the water from Bangor. It has a population of 3,376.
There are many small islands near the town, including Church Island. The Menai Heritage Bridges Exhibition celebrates the Menai Suspension Bridge, built by Thomas Telford, and the Britannia Bridge, built by Robert Stephenson.
Description and attractions
At the eastern edge of the town is Cwm Cadnant Dingle which is now by-passed by a modern bridge constructed in the 1970s. The Afon Cadnant drains into the Menai Strait at this point and this small estuary provides a natural haven for small boats crossing from the mainland. This was the location of the landing stage for the Bishops of Bangor who had their residence at Glyn Garth on Anglesey but whose cathedral was in Bangor on the mainland.
There are a number of small islands in the Menai Strait some of which are connected to the town by causeways, including Ynys Faelog, Ynys Gaint, Ynys Castell and Ynys y Bîg east of the suspension bridge and Church Island (Ynys Tysilio in Welsh) west of the bridge. The Isle of Anglesey Coastal Path passes along the waterfront.
Menai Bridge has several churches and chapels, including an English and Welsh Presbyterian church and a Catholic church. The town also has a primary school, Ysgol y Borth, and a large bilingual comprehensive school, Ysgol David Hughes.
Menai Bridge is home to the School of Ocean Sciences, part of Bangor University. Their research ship, the Prince Madog, is based at the pier when not at sea.
Attractions in Menai Bridge include the 14th-century Church of St Tysilio, St George's Pier, a butterfly house, Pili Palas, and the Plas Cadnant Hidden Gardens, a 200-acre (80 hectare) estate originally developed as a picturesque garden in the 1800s. The garden had been the site of restoration for twenty years. In December 2015, heavy rains caused flooding which washed away rare plants representing twenty years of work by Anthony Tavernor. Tavernor received some help to restore the garden, enabling him and his small staff to begin rebuilding and replanting the garden. The garden was able to reopen by Easter, 2016.
Listed buildings
There are over 30 buildings listed by Cadw of being of special importance. These include the suspension bridge itself, St. Mary's Church, the church of St. Tysilio, the Victoria Hotel, and the War memorial on Church Island and several individual houses and buildings
Glyn Garth
Menai Bridge includes the development along Beaumaris Road known as Glyn Garth. This was a favoured location for holiday houses for the wealthy from the Manchester and Liverpool areas in the late 19th century, and many large houses of that period remain. This was also where the Bishop of Bangor had his palace. The palace was demolished in the early 1960s and replaced by a block of flats, Glyn Garth Court, completed in 1966.
History
The town existed as Porthaethwy for centuries and still has a house which dates from the 17th century. The name derives from Porth (harbour) + Daethwy (the name of a local Celtic tribe and later of a local medieval commote). It is likely that a community existed here in Roman times as it is the shortest crossing of the Menai Strait.
In the 9th century, St Tysilio lived here as a hermit on Church Island.
A ferry across the Menai was first recorded in 1292. When the bridge opened in 1826, the ferry closed, but connections with the sea remained through the import, export and shipbuilding trades.
Lewis Carroll's Through the Looking Glass (1872) mentions the Menai Bridge in chapter 8 in a nonsense song.
From 1877 to 1920, the ship HMS Clio was docked at Menai Bridge; it was lent to the North Wales Society to teach young men the ways of seafaring. Many local people believed the ship was used for some type of prison, but this was not entirely true. The ship was home to young men who were in need of discipline to keep them from getting into serious trouble; some were sent to the Clio against their will. The young men on the Clio were not permitted to leave the ship; some of the corporal punishment administered was cruel. Stories about life on the Clio were commonplace among the residents of Menai Bridge; for many years, some mothers threatened their misbehaving children with being sent to live on the ship.
On 12 November 1918, Major Thomas Elmhirst (later Air Marshal Sir Thomas Elmhirst), commanding officer of RNAS Anglesey, flew airship SSZ73 under the Menai Bridge following the armistice at the end of World War I.
Carreg yr Halen
Carreg yr Halen is a small tidal island in the Menai Strait. Its centre lies approximately 20 metres offshore from the Belgian Promenade just upstream of the suspension bridge. Only the rocky tip of the island is visible at high spring tide but at low tide area of rock, sand and some seaweed are exposed which provides feeding ground for a variety of wading birds including oystercatcher, redshank, and curlew.
It is the site of one of the many ferry crossings of the Menai Strait which were in use prior to the construction of the suspension bridge in 1826
In 1914, Belgian refugees from Mechelen, who had settled in the area, built a promenade (the Belgian promenade) out of gratitude for the town's hospitality. The promenade was built along the Menai Strait from Ynys Tysilio (Church Island) to Carreg yr Halen and was completed in 1916. It was rebuilt in 1963. The ceremonial reopening in 1965 was performed by the only surviving refugee, Eduard Wilhelms. Most of the refugees lived at three houses in Menai Bridge, with 12 housed at the Village Hall in Llandegfan. Most of the men were skilled in marquetry.
A special celebration was held in 2014 at Menai Bridge to celebrate to centenary of the construction of the promenade.
TV location
Welsh-language production company, Rondo Media, has converted a disused garage into a fake row of shops in the centre of Menai Bridge as a film set for the soap opera Rownd a Rownd, shown on the Welsh-language channel S4C. They also film the show in schools in the town, Ysgol y Borth, and around the town itself.
Fair
The large car park to the north of the High Street is the "fair field". This is a piece of common land set aside for the holding of an annual fair called Ffair Borth, a tradition dating back to 1691. It started as a horse fair, and livestock trading was carried out until the 1970s. It was also a hiring fair. It was one of the year's great occasions for the folk of Anglesey and Arfon. The fair now features traditional fair rides. It comes to Menai Bridge on 24 October every year, unless it falls on a Sunday, in which case it is held on either 23 or 25 October. The fair stalls also take over most of the roads and streets in the town, making passage through the town very difficult.
A traditional verse goes:
Governance
There are two tiers of local government covering Menai Bridge, at community (town) and county level: Menai Bridge Town Council and Isle of Anglesey County Council. The town council is based at Canolfan Coed Cyrnol on Mona Road.
Administrative history
The community of Menai Bridge corresponds to the ancient parish of Llandysilio, which had its parish church at St Tysilio's Church on Church Island.
In 1884 the parish was made a local government district, administered by an elected local board. Although its area was defined as the parish of Llandysilio, the local government district took the name Menai Bridge. The local board also took over the functions of the Llandysilio improvement commissioners, which had been established in 1879 to manage certain areas of former common land in the parish. Such local government districts were reconstituted as urban districts under the Local Government Act 1894.
Menai Bridge Urban District was abolished in 1974, with its area instead becoming a community called Menai Bridge. District-level functions passed to Ynys Môn-Isle of Anglesey Borough Council, which in 1996 was reconstituted as a county council.
Notable people
Hugh Williams (1843–1911), a Welsh church historian, college tutor and Presbyterian minister.
Gwendolen Mason (1883-1977), Welsh harpist
Jane Helen Rowlands (1891–1955), scholar and missionary
Alun Owen (1925–94), screenwriter who wrote the script for The Beatles' A Hard Day's Night
| Technology | Bridges | null |
620358 | https://en.wikipedia.org/wiki/Iapetus%20Ocean | Iapetus Ocean | The Iapetus Ocean (; ) existed in the late Neoproterozoic and early Paleozoic eras of the geologic timescale (between 600 and 400 million years ago). It was in the southern hemisphere, between the paleocontinents of Laurentia, Baltica and Avalonia. The ocean disappeared with the Acadian, Caledonian and Taconic orogenies, when these three continents joined to form one big landmass called Euramerica. The "southern" Iapetus Ocean has been proposed to have closed with the Famatinian and Taconic orogenies, meaning a collision between Western Gondwana and Laurentia.
Because the Iapetus Ocean was positioned between continental masses that would at a much later time roughly form the opposite shores of the Atlantic Ocean, it can be seen as a sort of precursor of the Atlantic, and the process by which it opened shares many similarities with that of the Atlantic's initial opening in the Jurassic. The Iapetus Ocean was therefore named for the titan Iapetus, who in Greek mythology was the father of Atlas, after whom the Atlantic Ocean was named.
Research history
At the start of the 20th century, American paleontologist Charles Walcott noticed differences in early Paleozoic benthic trilobites of Laurentia (such as Olenellidae, the so-called "Pacific fauna"), as found in Scotland and western Newfoundland, and those of Baltica (such as Paradoxididae, often called the "Atlantic fauna"), as found in the southern parts of the British Isles and eastern Newfoundland. Geologists of the early 20th century presumed that a large trough, a so-called geosyncline, had existed between Scotland and England in the early Paleozoic, keeping the two sides separated.
With the development of plate tectonics in the 1960s, geologists such as Arthur Holmes and John Tuzo Wilson concluded that the Atlantic Ocean must have had a precursor before the time of Pangaea. Wilson also noticed that the Atlantic had opened at roughly the same place where its precursor ocean had closed. This led him to his Wilson cycle hypothesis.
Geodynamic history
Neoproterozoic origin
In many spots in Scandinavia basaltic dikes are found with ages between 670 and 650 million years. These are interpreted as evidence that by that time, rifting had started that would form the Iapetus Ocean. In Newfoundland and Labrador, the Long Range dikes are also thought to have formed during the formation of the Iapetus Ocean. It has been proposed that both the Fen Complex in Norway and the Alnö Complex in Sweden formed as consequence to mild extensional tectonics in the ancient continent of Baltica that followed the opening of the Iapetus Ocean.
The eastern Iapetus Ocean is believed to have opened around 590 Ma with the emplacement of the Central Iapetus Magmatic Province between Laurentia and Baltica. The southern Iapetus Ocean opened between Laurentia and southwestern Gondwana (now South America) about 550 Ma, close to the end of the Ediacaran period. At the time it did so the Adamastor Ocean further east closed. The opening of the Iapetus Ocean probably postdates the opening of the Puncoviscana Ocean, which is believed to have opened around 700 Ma as Laurentia drifted away from Amazonia, with the Iapetus Ocean being separated from the Puncoviscana Ocean by the ribbon-shaped Arequipa-Antofalla terrane. However, the formation of both oceans seems unrelated.
Paleozoic
Southwest of the Iapetus, a volcanic island arc evolved from the early Cambrian (540 million years ago) onward. This volcanic arc was formed above a subduction zone where the oceanic lithosphere of the Iapetus Ocean subducted southward under other oceanic lithosphere. From Cambrian times (about 550 million years ago) the western Iapetus Ocean began to grow progressively narrower due to this subduction. The same happened further north and east, where Avalonia and Baltica began to move towards Laurentia from the Ordovician (488–444 million years ago) onward.
Trilobite faunas of the continental shelves of Baltica and Laurentia are still very different in the Ordovician, but Silurian faunas show progressive mixing of species from both sides, because the continents moved closer together.
In the west, the Iapetus Ocean closed with the Taconic orogeny (480-430 million years ago), when the volcanic island arc collided with Laurentia. Some authors consider the oceanic basin south of the island arc also a part of the Iapetus, this branch closed during the later Acadian orogeny, when Avalonia collided with Laurentia.
It has been suggested that the southern Iapetus Ocean closed during a continental collision between Laurentia and Western Gondwana (South America). If factual the Taconic orogen would be the northward continuation of the Famatinian orogen exposed in Argentina.
Meanwhile, the eastern parts had closed too: the Tornquist Sea between Avalonia and Baltica already during the late Ordovician, the main branch between Baltica-Avalonia and Laurentia during the Grampian and Scandian phases of the Caledonian orogeny (440–420 million years ago).
At the end of the Silurian period (c. 420 million years ago) the Iapetus Ocean had completely disappeared and the combined mass of the three continents formed the "new" continent of Laurasia, which would itself be the northern component of the singular supercontinent of Pangaea.
| Physical sciences | Paleogeography | Earth science |
620629 | https://en.wikipedia.org/wiki/Eoraptor | Eoraptor | Eoraptor () is a genus of small, lightly built, basal sauropodomorph dinosaur. One of the earliest-known dinosaurs and one of the earliest sauropodomorphs, it lived approximately 231 to 228 million years ago, during the Late Triassic in Western Gondwana, in the region that is now northwestern Argentina. The type and only species, Eoraptor lunensis, was first described in 1993, and is known from an almost complete and well-preserved skeleton and several fragmentary ones. Eoraptor had multiple tooth shapes, which suggests that it was omnivorous.
History of discovery
The bones of this primitive dinosaur were first discovered in 1991, by University of San Juan paleontologist Ricardo Martínez, during field work conducted by the University of Chicago and the University of San Juan. The holotype specimen PVSJ 512 was discovered in muddy siltstone belonging to the Cancha de Bochas Member of the Ischigualasto Formation in Argentina. The fossils in this formation were deposited in the Carnian stage of the Triassic period, approximately 235 to 228 million years ago. It took almost 12 months to collect the holotype, which was then shipped to the Field Museum of Natural History in Chicago for preparation by William F. Simpson and Bob Masek. The fossil was first put on display in Chicago and was then returned to San Juan, Argentina, where it went on display at the Museum of Natural Sciences.
The genus Eoraptor was described and named by Paul Sereno, Catherine Forster, Raymond R. Rogers, and Alfredo M. Monetta in 1993. The name is derived from the Greek word () meaning 'dawn', a reference to its primitive nature, and the Latin word meaning 'plunderer', a reference to its presumed carnivorous nature and its grasping hand. The specific name lunensis is derived from the Latin words ('moon') and the suffix ('inhabitant'), a reference to its place of discovery: the ('Valley of the Moon'), so named for its arid, otherworldly appearance evocative of a lunar landscape. The type species Eoraptor lunensis means 'dawn plunderer from the Valley of the Moon'.
Description
Eoraptor was a small dinosaur, with the known specimens measuring in length, and weighing around or less than . It had a lightly built skull with a slightly enlarged external naris. As in early sauropodomorphs such as Buriolestes and Pampadromaeus and coelophysoids (which would appear millions of years later), Eoraptor had a kink in its upper jaws, between the maxilla and the premaxilla. Paul Sereno et al. (2013) observed that the lower jaw had a mid-mandibular joint. It ran digitigrade, and upright on its hind legs. The femur of the holotype specimen PVSJ 512 is , and the tibia is , suggesting that it was a fast runner. Its forelimbs are only half the length of its hindlimbs, suggesting that it was bipedal. All of its long bones have hollow shafts. Eoraptor had five digits on each 'hand', the three longest of which ended in large claws and were presumably used to handle prey. Scientists have surmised that the fourth and fifth digits were too tiny to be of any use in hunting. The ilium is supported by three sacral vertebrae (atypical of the plesiomorphic two sacrals of basal sauropodomorphs), unlike that of the coeval Herrerasaurus which is supported by only two sacrals, a basal trait. Eoraptor had vertebral centra that are hollow, a feature present in some of its ancestors.
The original describers, Paul Sereno et al. (1993), supported the notion that Eoraptor was an adult specimen based on the closure of sutures in the vertebral column, and the partial fusion of the scapulocoracoid. Bonaparte (1996) interpreted the relatively large orbital opening in the skull as a juvenile trait. Ronald Tykoski agreed (2005) and suggested that certain skull features of the type specimen suggested that it was young, specifically, the skull bones are not completely fused, relatively large orbits, and a short snout. Later Sereno et al. (2013) considered the type specimen as a young adult approaching skeletal maturity, considering that it contained traits of both maturity and immaturity.
According to Sereno et al. (1993), Eoraptor can be distinguished based on the fact that its premaxillary and anterior maxillary teeth are leaf-shaped, the external nares are slightly enlarged, and the premaxilla is observed to have a slender posterolateral process. Max Langer and Michael Benton (2006) noted that Eoraptor can be distinguished based on the fact that the proximal part of its fibula is extremely transversely compressed.
Classification
In 1993 Paul Sereno and his colleagues described and named Eoraptor, and determined it to be one of the earliest dinosaurs. Its age was determined by several factors, not least because it lacked the specialized features of any of the major groups of later dinosaurs, including its lack of specialized predatory features. In 1995, Sereno posited that Eoraptor is the earliest-recorded theropod, and is closest to "the hypothetical dinosaurian condition than any other dinosaurian subgroup." The precise placement of Eoraptor within Dinosauria has been unstable, with opinion often varying between a basal saurischian and a basal theropod. When it was first described by Sereno and Forster in 1993, it was regarded as a theropod, based on its "functionally tridactyl hand" and other anatomical features. In 2011, a study conducted by Hans-Dieter Sues, Sterling J. Nesbitt, David S. Berman and Amy C. Henrici featuring a description of Daemonosaurus, also concluded that there is now enough fossil evidence to confidently classify Eoraptor as a theropod. The study noted that the "transitional suite of character states" of the recently discovered dinosaurs, Daemonosaurus and Tawa further support that Eoraptor is a basal theropod, and not a basal saurischian or a basal sauropodomorph. On the other hand, several studies from 2012 onward have recovered Eoraptor as an early sauropodomorph, rather than a theropod. The following phylogenetic tree illustrates the relationships of Eoraptor among the major theropod groups based on various studies conducted in the 2010s.
Philip Currie (1997) found Eoraptor anatomically closer to what would be considered the ancestral morphotype of both saurischian and ornithischian dinosaurs. In 2011, Martinez et al. (the team that described Eodromaeus) found Eoraptor to be a basal sauropodomorph, with characteristic features from the group. Michael Benton expressed his hesitation to this, and claimed that it is "quite a shift" to remove Eoraptor from Theropoda and then place it in Sauropodomorpha. A subsequent study by Apaldetti, Martinez, Alcober, and Pol published in 2011 found Eoraptor to be a saurischian close to sauropodomorphs and theropods, though was unable to resolve which of the two branches, if either, it fell within. Sereno et al. (2013) redescribed the holotype skeleton and concluded that Eoraptor was not a theropod but a basal sauropodomorph, consistent with the earlier observation made by Martinez et al. (2011).
A large phylogenetic analysis of early dinosaurs by Matthew Baron, David Norman and Paul Barrett (2017) found Eoraptor to be the earliest diverging member of Theropoda, within the larger clade Ornithoscelida. A phylogenetic analysis published with the description of new Buriolestes remains in 2018, based on Langer et al. (2017) placed Eoraptor in a clade of early sauropodomorphs, alongside Buriolestes, Panphagia, Pampadromaeus, and Saturnalia.
Paleobiology
Eoraptor is thought to have been an omnivore, although its dentition is quite similar to that of Buriolestes, which is considered carnivorous. It was a swift sprinter and, upon catching its prey, it would use claws and teeth to tear the prey apart. Unlike later, carnivorous dinosaurs, it lacked a sliding joint at the articulation of the lower jaw, with which to hold large prey. Furthermore, only some of its teeth were curved and saw-edged, unlike those in the mouths of later theropods. The heterodont dentition of Eoraptor consists of both serrated, recurved teeth in the upper jaw, like the teeth of theropods, and leaf-shaped teeth in the lower jaw, like the teeth of basal sauropodomorphs. Eoraptor had 4 teeth in the premaxilla and 18 teeth in the maxilla, a dental formula not dissimilar to that of Herrerasaurus.
Paleoecology
During the Late Triassic period, the Ischigualasto Formation was a volcanically active floodplain covered by forests, with a warm and humid climate, but subject to seasonal variations including strong rainfall. Vegetation consisted of ferns, horsetails, and giant conifers, which formed highland forests along the banks of rivers. Herrerasaurus remains appear to have been the most common among the carnivores of the Ischigualasto Formation. Sereno (1993) noted that Eoraptor was found in "close association" with therapsids, rauisuchians, archosaurs, Saurosuchus and the dinosaurs Herrerasaurus and Pisanosaurus, all of whom lived in its paleoenvironment. Herbivores were represented by rhynchosaurs such as Hyperodapedon; aetosaurs; cynodonts like Chiniquodon, kannemeyeriid dicynodonts such as Ischigualastia; and traversodontids such as Exaeretodon. These non-dinosaurian herbivores were much more abundant than early dinosaurs. Dinosaur fossils, including those of Eoraptor only represent approximately 6% of the total sample that has been recovered from the Ischigualasto Formation (Rogers et al., 1993), which suggests that dinosaurs were less numerous than other tetrapods.
| Biology and health sciences | Dinosaurs | Animals |
620667 | https://en.wikipedia.org/wiki/Methodology | Methodology | In its most common sense, methodology is the study of research methods. However, the term can also refer to the methods themselves or to the philosophical discussion of associated background assumptions. A method is a structured procedure for bringing about a certain goal, like acquiring knowledge or verifying knowledge claims. This normally involves various steps, like choosing a sample, collecting data from this sample, and interpreting the data. The study of methods concerns a detailed description and analysis of these processes. It includes evaluative aspects by comparing different methods. This way, it is assessed what advantages and disadvantages they have and for what research goals they may be used. These descriptions and evaluations depend on philosophical background assumptions. Examples are how to conceptualize the studied phenomena and what constitutes evidence for or against them. When understood in the widest sense, methodology also includes the discussion of these more abstract issues.
Methodologies are traditionally divided into quantitative and qualitative research. Quantitative research is the main methodology of the natural sciences. It uses precise numerical measurements. Its goal is usually to find universal laws used to make predictions about future events. The dominant methodology in the natural sciences is called the scientific method. It includes steps like observation and the formulation of a hypothesis. Further steps are to test the hypothesis using an experiment, to compare the measurements to the expected results, and to publish the findings.
Qualitative research is more characteristic of the social sciences and gives less prominence to exact numerical measurements. It aims more at an in-depth understanding of the meaning of the studied phenomena and less at universal and predictive laws. Common methods found in the social sciences are surveys, interviews, focus groups, and the nominal group technique. They differ from each other concerning their sample size, the types of questions asked, and the general setting. In recent decades, many social scientists have started using mixed-methods research, which combines quantitative and qualitative methodologies.
Many discussions in methodology concern the question of whether the quantitative approach is superior, especially whether it is adequate when applied to the social domain. A few theorists reject methodology as a discipline in general. For example, some argue that it is useless since methods should be used rather than studied. Others hold that it is harmful because it restricts the freedom and creativity of researchers. Methodologists often respond to these objections by claiming that a good methodology helps researchers arrive at reliable theories in an efficient way. The choice of method often matters since the same factual material can lead to different conclusions depending on one's method. Interest in methodology has risen in the 20th century due to the increased importance of interdisciplinary work and the obstacles hindering efficient cooperation.
Definitions
The term "methodology" is associated with a variety of meanings. In its most common usage, it refers either to a method, to the field of inquiry studying methods, or to philosophical discussions of background assumptions involved in these processes. Some researchers distinguish methods from methodologies by holding that methods are modes of data collection while methodologies are more general research strategies that determine how to conduct a research project. In this sense, methodologies include various theoretical commitments about the intended outcomes of the investigation.
As method
The term "methodology" is sometimes used as a synonym for the term "method". A method is a way of reaching some predefined goal. It is a planned and structured procedure for solving a theoretical or practical problem. In this regard, methods stand in contrast to free and unstructured approaches to problem-solving. For example, descriptive statistics is a method of data analysis, radiocarbon dating is a method of determining the age of organic objects, sautéing is a method of cooking, and project-based learning is an educational method. The term "technique" is often used as a synonym both in the academic and the everyday discourse. Methods usually involve a clearly defined series of decisions and actions to be used under certain circumstances, usually expressable as a sequence of repeatable instructions. The goal of following the steps of a method is to bring about the result promised by it. In the context of inquiry, methods may be defined as systems of rules and procedures to discover regularities of nature, society, and thought. In this sense, methodology can refer to procedures used to arrive at new knowledge or to techniques of verifying and falsifying pre-existing knowledge claims. This encompasses various issues pertaining both to the collection of data and their analysis. Concerning the collection, it involves the problem of sampling and of how to go about the data collection itself, like surveys, interviews, or observation. There are also numerous methods of how the collected data can be analyzed using statistics or other ways of interpreting it to extract interesting conclusions.
As study of methods
However, many theorists emphasize the differences between the terms "method" and "methodology". In this regard, methodology may be defined as "the study or description of methods" or as "the analysis of the principles of methods, rules, and postulates employed by a discipline". This study or analysis involves uncovering assumptions and practices associated with the different methods and a detailed description of research designs and hypothesis testing. It also includes evaluative aspects: forms of data collection, measurement strategies, and ways to analyze data are compared and their advantages and disadvantages relative to different research goals and situations are assessed. In this regard, methodology provides the skills, knowledge, and practical guidance needed to conduct scientific research in an efficient manner. It acts as a guideline for various decisions researchers need to take in the scientific process.
Methodology can be understood as the middle ground between concrete particular methods and the abstract and general issues discussed by the philosophy of science. In this regard, methodology comes after formulating a research question and helps the researchers decide what methods to use in the process. For example, methodology should assist the researcher in deciding why one method of sampling is preferable to another in a particular case or which form of data analysis is likely to bring the best results. Methodology achieves this by explaining, evaluating and justifying methods. Just as there are different methods, there are also different methodologies. Different methodologies provide different approaches to how methods are evaluated and explained and may thus make different suggestions on what method to use in a particular case.
According to Aleksandr Georgievich Spirkin, "[a] methodology is a system of principles and general ways of organising and structuring theoretical and practical activity, and also the theory of this system". Helen Kara defines methodology as "a contextual framework for research, a coherent and logical scheme based on views, beliefs, and values, that guides the choices researchers make". Ginny E. Garcia and Dudley L. Poston understand methodology either as a complex body of rules and postulates guiding research or as the analysis of such rules and procedures. As a body of rules and postulates, a methodology defines the subject of analysis as well as the conceptual tools used by the analysis and the limits of the analysis. Research projects are usually governed by a structured procedure known as the research process. The goal of this process is given by a research question, which determines what kind of information one intends to acquire.
As discussion of background assumptions
Some theorists prefer an even wider understanding of methodology that involves not just the description, comparison, and evaluation of methods but includes additionally more general philosophical issues. One reason for this wider approach is that discussions of when to use which method often take various background assumptions for granted, for example, concerning the goal and nature of research. These assumptions can at times play an important role concerning which method to choose and how to follow it. For example, Thomas Kuhn argues in his The Structure of Scientific Revolutions that sciences operate within a framework or a paradigm that determines which questions are asked and what counts as good science. This concerns philosophical disagreements both about how to conceptualize the phenomena studied, what constitutes evidence for and against them, and what the general goal of researching them is. So in this wider sense, methodology overlaps with philosophy by making these assumptions explicit and presenting arguments for and against them. According to C. S. Herrman, a good methodology clarifies the structure of the data to be analyzed and helps the researchers see the phenomena in a new light. In this regard, a methodology is similar to a paradigm. A similar view is defended by Spirkin, who holds that a central aspect of every methodology is the world view that comes with it.
The discussion of background assumptions can include metaphysical and ontological issues in cases where they have important implications for the proper research methodology. For example, a realist perspective considering the observed phenomena as an external and independent reality is often associated with an emphasis on empirical data collection and a more distanced and objective attitude. Idealists, on the other hand, hold that external reality is not fully independent of the mind and tend, therefore, to include more subjective tendencies in the research process as well.
For the quantitative approach, philosophical debates in methodology include the distinction between the inductive and the hypothetico-deductive interpretation of the scientific method. For qualitative research, many basic assumptions are tied to philosophical positions such as hermeneutics, pragmatism, Marxism, critical theory, and postmodernism. According to Kuhn, an important factor in such debates is that the different paradigms are incommensurable. This means that there is no overarching framework to assess the conflicting theoretical and methodological assumptions. This critique puts into question various presumptions of the quantitative approach associated with scientific progress based on the steady accumulation of data.
Other discussions of abstract theoretical issues in the philosophy of science are also sometimes included. This can involve questions like how and whether scientific research differs from fictional writing as well as whether research studies objective facts rather than constructing the phenomena it claims to study. In the latter sense, some methodologists have even claimed that the goal of science is less to represent a pre-existing reality and more to bring about some kind of social change in favor of repressed groups in society.
Related terms and issues
Viknesh Andiappan and Yoke Kin Wan use the field of process systems engineering to distinguish the term "methodology" from the closely related terms "approach", "method", "procedure", and "technique". On their view, "approach" is the most general term. It can be defined as "a way or direction used to address a problem based on a set of assumptions". An example is the difference between hierarchical approaches, which consider one task at a time in a hierarchical manner, and concurrent approaches, which consider them all simultaneously. Methodologies are a little more specific. They are general strategies needed to realize an approach and may be understood as guidelines for how to make choices. Often the term "framework" is used as a synonym. A method is a still more specific way of practically implementing the approach. Methodologies provide the guidelines that help researchers decide which method to follow. The method itself may be understood as a sequence of techniques. A technique is a step taken that can be observed and measured. Each technique has some immediate result. The whole sequence of steps is termed a "procedure". A similar but less complex characterization is sometimes found in the field of language teaching, where the teaching process may be described through a three-level conceptualization based on "approach", "method", and "technique".
One question concerning the definition of methodology is whether it should be understood as a descriptive or a normative discipline. The key difference in this regard is whether methodology just provides a value-neutral description of methods or what scientists actually do. Many methodologists practice their craft in a normative sense, meaning that they express clear opinions about the advantages and disadvantages of different methods. In this regard, methodology is not just about what researchers actually do but about what they ought to do or how to perform good research.
Types
Theorists often distinguish various general types or approaches to methodology. The most influential classification contrasts quantitative and qualitative methodology.
Quantitative and qualitative
Quantitative research is closely associated with the natural sciences. It is based on precise numerical measurements, which are then used to arrive at exact general laws. This precision is also reflected in the goal of making predictions that can later be verified by other researchers. Examples of quantitative research include physicists at the Large Hadron Collider measuring the mass of newly created particles and positive psychologists conducting an online survey to determine the correlation between income and self-assessed well-being.
Qualitative research is characterized in various ways in the academic literature but there are very few precise definitions of the term. It is often used in contrast to quantitative research for forms of study that do not quantify their subject matter numerically. However, the distinction between these two types is not always obvious and various theorists have argued that it should be understood as a continuum and not as a dichotomy. A lot of qualitative research is concerned with some form of human experience or behavior, in which case it tends to focus on a few individuals and their in-depth understanding of the meaning of the studied phenomena. Examples of the qualitative method are a market researcher conducting a focus group in order to learn how people react to a new product or a medical researcher performing an unstructured in-depth interview with a participant from a new experimental therapy to assess its potential benefits and drawbacks. It is also used to improve quantitative research, such as informing data collection materials and questionnaire design. Qualitative research is frequently employed in fields where the pre-existing knowledge is inadequate. This way, it is possible to get a first impression of the field and potential theories, thus paving the way for investigating the issue in further studies.
Quantitative methods dominate in the natural sciences but both methodologies are used in the social sciences. Some social scientists focus mostly on one method while others try to investigate the same phenomenon using a variety of different methods. It is central to both approaches how the group of individuals used for the data collection is selected. This process is known as sampling. It involves the selection of a subset of individuals or phenomena to be measured. Important in this regard is that the selected samples are representative of the whole population, i.e. that no significant biases were involved when choosing. If this is not the case, the data collected does not reflect what the population as a whole is like. This affects generalizations and predictions drawn from the biased data. The number of individuals selected is called the sample size. For qualitative research, the sample size is usually rather small, while quantitative research tends to focus on big groups and collecting a lot of data. After the collection, the data needs to be analyzed and interpreted to arrive at interesting conclusions that pertain directly to the research question. This way, the wealth of information obtained is summarized and thus made more accessible to others. Especially in the case of quantitative research, this often involves the application of some form of statistics to make sense of the numerous individual measurements.
Many discussions in the history of methodology center around the quantitative methods used by the natural sciences. A central question in this regard is to what extent they can be applied to other fields, like the social sciences and history. The success of the natural sciences was often seen as an indication of the superiority of the quantitative methodology and used as an argument to apply this approach to other fields as well. However, this outlook has been put into question in the more recent methodological discourse. In this regard, it is often argued that the paradigm of the natural sciences is a one-sided development of reason, which is not equally well suited to all areas of inquiry. The divide between quantitative and qualitative methods in the social sciences is one consequence of this criticism.
Which method is more appropriate often depends on the goal of the research. For example, quantitative methods usually excel for evaluating preconceived hypotheses that can be clearly formulated and measured. Qualitative methods, on the other hand, can be used to study complex individual issues, often with the goal of formulating new hypotheses. This is especially relevant when the existing knowledge of the subject is inadequate. Important advantages of quantitative methods include precision and reliability. However, they have often difficulties in studying very complex phenomena that are commonly of interest to the social sciences. Additional problems can arise when the data is misinterpreted to defend conclusions that are not directly supported by the measurements themselves. In recent decades, many researchers in the social sciences have started combining both methodologies. This is known as mixed-methods research. A central motivation for this is that the two approaches can complement each other in various ways: some issues are ignored or too difficult to study with one methodology and are better approached with the other. In other cases, both approaches are applied to the same issue to produce more comprehensive and well-rounded results.
Qualitative and quantitative research are often associated with different research paradigms and background assumptions. Qualitative researchers often use an interpretive or critical approach while quantitative researchers tend to prefer a positivistic approach. Important disagreements between these approaches concern the role of objectivity and hard empirical data as well as the research goal of predictive success rather than in-depth understanding or social change.
Others
Various other classifications have been proposed. One distinguishes between substantive and formal methodologies. Substantive methodologies tend to focus on one specific area of inquiry. The findings are initially restricted to this specific field but may be transferrable to other areas of inquiry. Formal methodologies, on the other hand, are based on a variety of studies and try to arrive at more general principles applying to different fields. They may also give particular prominence to the analysis of the language of science and the formal structure of scientific explanation. A closely related classification distinguishes between philosophical, general scientific, and special scientific methods.
One type of methodological outlook is called "proceduralism". According to it, the goal of methodology is to boil down the research process to a simple set of rules or a recipe that automatically leads to good research if followed precisely. However, it has been argued that, while this ideal may be acceptable for some forms of quantitative research, it fails for qualitative research. One argument for this position is based on the claim that research is not a technique but a craft that cannot be achieved by blindly following a method. In this regard, research depends on forms of creativity and improvisation to amount to good science.
Other types include inductive, deductive, and transcendental methods. Inductive methods are common in the empirical sciences and proceed through inductive reasoning from many particular observations to arrive at general conclusions, often in the form of universal laws. Deductive methods, also referred to as axiomatic methods, are often found in formal sciences, such as geometry. They start from a set of self-evident axioms or first principles and use deduction to infer interesting conclusions from these axioms. Transcendental methods are common in Kantian and post-Kantian philosophy. They start with certain particular observations. It is then argued that the observed phenomena can only exist if their conditions of possibility are fulfilled. This way, the researcher may draw general psychological or metaphysical conclusions based on the claim that the phenomenon would not be observable otherwise.
Importance
It has been argued that a proper understanding of methodology is important for various issues in the field of research. They include both the problem of conducting efficient and reliable research as well as being able to validate knowledge claims by others. Method is often seen as one of the main factors of scientific progress. This is especially true for the natural sciences where the developments of experimental methods in the 16th and 17th century are often seen as the driving force behind the success and prominence of the natural sciences. In some cases, the choice of methodology may have a severe impact on a research project. The reason is that very different and sometimes even opposite conclusions may follow from the same factual material based on the chosen methodology.
Aleksandr Georgievich Spirkin argues that methodology, when understood in a wide sense, is of great importance since the world presents us with innumerable entities and relations between them. Methods are needed to simplify this complexity and find a way of mastering it. On the theoretical side, this concerns ways of forming true beliefs and solving problems. On the practical side, this concerns skills of influencing nature and dealing with each other. These different methods are usually passed down from one generation to the next. Spirkin holds that the interest in methodology on a more abstract level arose in attempts to formalize these techniques to improve them as well as to make it easier to use them and pass them on. In the field of research, for example, the goal of this process is to find reliable means to acquire knowledge in contrast to mere opinions acquired by unreliable means. In this regard, "methodology is a way of obtaining and building up ... knowledge".
Various theorists have observed that the interest in methodology has risen significantly in the 20th century. This increased interest is reflected not just in academic publications on the subject but also in the institutionalized establishment of training programs focusing specifically on methodology. This phenomenon can be interpreted in different ways. Some see it as a positive indication of the topic's theoretical and practical importance. Others interpret this interest in methodology as an excessive preoccupation that draws time and energy away from doing research on concrete subjects by applying the methods instead of researching them. This ambiguous attitude towards methodology is sometimes even exemplified in the same person. Max Weber, for example, criticized the focus on methodology during his time while making significant contributions to it himself. Spirkin believes that one important reason for this development is that contemporary society faces many global problems. These problems cannot be solved by a single researcher or a single discipline but are in need of collaborative efforts from many fields. Such interdisciplinary undertakings profit a lot from methodological advances, both concerning the ability to understand the methods of the respective fields and in relation to developing more homogeneous methods equally used by all of them.
Criticism
Most criticism of methodology is directed at one specific form or understanding of it. In such cases, one particular methodological theory is rejected but not methodology at large when understood as a field of research comprising many different theories. In this regard, many objections to methodology focus on the quantitative approach, specifically when it is treated as the only viable approach. Nonetheless, there are also more fundamental criticisms of methodology in general. They are often based on the idea that there is little value to abstract discussions of methods and the reasons cited for and against them. In this regard, it may be argued that what matters is the correct employment of methods and not their meticulous study. Sigmund Freud, for example, compared methodologists to "people who clean their glasses so thoroughly that they never have time to look through them". According to C. Wright Mills, the practice of methodology often degenerates into a "fetishism of method and technique".
Some even hold that methodological reflection is not just a waste of time but actually has negative side effects. Such an argument may be defended by analogy to other skills that work best when the agent focuses only on employing them. In this regard, reflection may interfere with the process and lead to avoidable mistakes. According to an example by Gilbert Ryle, "[w]e run, as a rule, worse, not better, if we think a lot about our feet". A less severe version of this criticism does not reject methodology per se but denies its importance and rejects an intense focus on it. In this regard, methodology has still a limited and subordinate utility but becomes a diversion or even counterproductive by hindering practice when given too much emphasis.
Another line of criticism concerns more the general and abstract nature of methodology. It states that the discussion of methods is only useful in concrete and particular cases but not concerning abstract guidelines governing many or all cases. Some anti-methodologists reject methodology based on the claim that researchers need freedom to do their work effectively. But this freedom may be constrained and stifled by "inflexible and inappropriate guidelines". For example, according to Kerry Chamberlain, a good interpretation needs creativity to be provocative and insightful, which is prohibited by a strictly codified approach. Chamberlain uses the neologism "methodolatry" to refer to this alleged overemphasis on methodology. Similar arguments are given in Paul Feyerabend's book "Against Method".
However, these criticisms of methodology in general are not always accepted. Many methodologists defend their craft by pointing out how the efficiency and reliability of research can be improved through a proper understanding of methodology.
A criticism of more specific forms of methodology is found in the works of the sociologist Howard S. Becker. He is quite critical of methodologists based on the claim that they usually act as advocates of one particular method usually associated with quantitative research. An often-cited quotation in this regard is that "[m]ethodology is too important to be left to methodologists". Alan Bryman has rejected this negative outlook on methodology. He holds that Becker's criticism can be avoided by understanding methodology as an inclusive inquiry into all kinds of methods and not as a mere doctrine for converting non-believers to one's preferred method.
In different fields
Part of the importance of methodology is reflected in the number of fields to which it is relevant. They include the natural sciences and the social sciences as well as philosophy and mathematics.
Natural sciences
The dominant methodology in the natural sciences (like astronomy, biology, chemistry, geoscience, and physics) is called the scientific method. Its main cognitive aim is usually seen as the creation of knowledge, but various closely related aims have also been proposed, like understanding, explanation, or predictive success. Strictly speaking, there is no one single scientific method. In this regard, the expression "scientific method" refers not to one specific procedure but to different general or abstract methodological aspects characteristic of all the aforementioned fields. Important features are that the problem is formulated in a clear manner and that the evidence presented for or against a theory is public, reliable, and replicable. The last point is important so that other researchers are able to repeat the experiments to confirm or disconfirm the initial study. For this reason, various factors and variables of the situation often have to be controlled to avoid distorting influences and to ensure that subsequent measurements by other researchers yield the same results. The scientific method is a quantitative approach that aims at obtaining numerical data. This data is often described using mathematical formulas. The goal is usually to arrive at some universal generalizations that apply not just to the artificial situation of the experiment but to the world at large. Some data can only be acquired using advanced measurement instruments. In cases where the data is very complex, it is often necessary to employ sophisticated statistical techniques to draw conclusions from it.
The scientific method is often broken down into several steps. In a typical case, the procedure starts with regular observation and the collection of information. These findings then lead the scientist to formulate a hypothesis describing and explaining the observed phenomena. The next step consists in conducting an experiment designed for this specific hypothesis. The actual results of the experiment are then compared to the expected results based on one's hypothesis. The findings may then be interpreted and published, either as a confirmation or disconfirmation of the initial hypothesis.
Two central aspects of the scientific method are observation and experimentation. This distinction is based on the idea that experimentation involves some form of manipulation or intervention. This way, the studied phenomena are actively created or shaped. For example, a biologist inserting viral DNA into a bacterium is engaged in a form of experimentation. Pure observation, on the other hand, involves studying independent entities in a passive manner. This is the case, for example, when astronomers observe the orbits of astronomical objects far away. Observation played the main role in ancient science. The scientific revolution in the 16th and 17th century affected a paradigm change that gave a much more central role to experimentation in the scientific methodology. This is sometimes expressed by stating that modern science actively "puts questions to nature". While the distinction is usually clear in the paradigmatic cases, there are also many intermediate cases where it is not obvious whether they should be characterized as observation or as experimentation.
A central discussion in this field concerns the distinction between the inductive and the hypothetico-deductive methodology. The core disagreement between these two approaches concerns their understanding of the confirmation of scientific theories. The inductive approach holds that a theory is confirmed or supported by all its positive instances, i.e. by all the observations that exemplify it. For example, the observations of many white swans confirm the universal hypothesis that "all swans are white". The hypothetico-deductive approach, on the other hand, focuses not on positive instances but on deductive consequences of the theory. This way, the researcher uses deduction before conducting an experiment to infer what observations they expect. These expectations are then compared to the observations they actually make. This approach often takes a negative form based on falsification. In this regard, positive instances do not confirm a hypothesis but negative instances disconfirm it. Positive indications that the hypothesis is true are only given indirectly if many attempts to find counterexamples have failed. A cornerstone of this approach is the null hypothesis, which assumes that there is no connection (see causality) between whatever is being observed. It is up to the researcher to do all they can to disprove their own hypothesis through relevant methods or techniques, documented in a clear and replicable process. If they fail to do so, it can be concluded that the null hypothesis is false, which provides support for their own hypothesis about the relation between the observed phenomena.
Social sciences
Significantly more methodological variety is found in the social sciences, where both quantitative and qualitative approaches are used. They employ various forms of data collection, such as surveys, interviews, focus groups, and the nominal group technique. Surveys belong to quantitative research and usually involve some form of questionnaire given to a large group of individuals. It is paramount that the questions are easily understandable by the participants since the answers might not have much value otherwise. Surveys normally restrict themselves to closed questions in order to avoid various problems that come with the interpretation of answers to open questions. They contrast in this regard to interviews, which put more emphasis on the individual participant and often involve open questions. Structured interviews are planned in advance and have a fixed set of questions given to each individual. They contrast with unstructured interviews, which are closer to a free-flow conversation and require more improvisation on the side of the interviewer for finding interesting and relevant questions. Semi-structured interviews constitute a middle ground: they include both predetermined questions and questions not planned in advance. Structured interviews make it easier to compare the responses of the different participants and to draw general conclusions. However, they also limit what may be discovered and thus constrain the investigation in many ways. Depending on the type and depth of the interview, this method belongs either to quantitative or to qualitative research. The terms research conversation and muddy interview have been used to describe interviews conducted in informal settings which may not occur purely for the purposes of data collection. Some researcher employ the go-along method by conducting interviews while they and the participants navigate through and engage with their environment.
Focus groups are a qualitative research method often used in market research. They constitute a form of group interview involving a small number of demographically similar people. Researchers can use this method to collect data based on the interactions and responses of the participants. The interview often starts by asking the participants about their opinions on the topic under investigation, which may, in turn, lead to a free exchange in which the group members express and discuss their personal views. An important advantage of focus groups is that they can provide insight into how ideas and understanding operate in a cultural context. However, it is usually difficult to use these insights to discern more general patterns true for a wider public. One advantage of focus groups is that they can help the researcher identify a wide range of distinct perspectives on the issue in a short time. The group interaction may also help clarify and expand interesting contributions. One disadvantage is due to the moderator's personality and group effects, which may influence the opinions stated by the participants. When applied to cross-cultural settings, cultural and linguistic adaptations and group composition considerations are important to encourage greater participation in the group discussion.
The nominal group technique is similar to focus groups with a few important differences. The group often consists of experts in the field in question. The group size is similar but the interaction between the participants is more structured. The goal is to determine how much agreement there is among the experts on the different issues. The initial responses are often given in written form by each participant without a prior conversation between them. In this manner, group effects potentially influencing the expressed opinions are minimized. In later steps, the different responses and comments may be discussed and compared to each other by the group as a whole.
Most of these forms of data collection involve some type of observation. Observation can take place either in a natural setting, i.e. the field, or in a controlled setting such as a laboratory. Controlled settings carry with them the risk of distorting the results due to their artificiality. Their advantage lies in precisely controlling the relevant factors, which can help make the observations more reliable and repeatable. Non-participatory observation involves a distanced or external approach. In this case, the researcher focuses on describing and recording the observed phenomena without causing or changing them, in contrast to participatory observation.
An important methodological debate in the field of social sciences concerns the question of whether they deal with hard, objective, and value-neutral facts, as the natural sciences do. Positivists agree with this characterization, in contrast to interpretive and critical perspectives on the social sciences. According to William Neumann, positivism can be defined as "an organized method for combining deductive logic with precise empirical observations of individual behavior in order to discover and confirm a set of probabilistic causal laws that can be used to predict general patterns of human activity". This view is rejected by interpretivists. Max Weber, for example, argues that the method of the natural sciences is inadequate for the social sciences. Instead, more importance is placed on meaning and how people create and maintain their social worlds. The critical methodology in social science is associated with Karl Marx and Sigmund Freud. It is based on the assumption that many of the phenomena studied using the other approaches are mere distortions or surface illusions. It seeks to uncover deeper structures of the material world hidden behind these distortions. This approach is often guided by the goal of helping people effect social changes and improvements.
Philosophy
Philosophical methodology is the metaphilosophical field of inquiry studying the methods used in philosophy. These methods structure how philosophers conduct their research, acquire knowledge, and select between competing theories. It concerns both descriptive issues of what methods have been used by philosophers in the past and normative issues of which methods should be used. Many philosophers emphasize that these methods differ significantly from the methods found in the natural sciences in that they usually do not rely on experimental data obtained through measuring equipment. Which method one follows can have wide implications for how philosophical theories are constructed, what theses are defended, and what arguments are cited in favor or against. In this regard, many philosophical disagreements have their source in methodological disagreements. Historically, the discovery of new methods, like methodological skepticism and the phenomenological method, has had important impacts on the philosophical discourse.
A great variety of methods has been employed throughout the history of philosophy. Methodological skepticism gives special importance to the role of systematic doubt. This way, philosophers try to discover absolutely certain first principles that are indubitable. The geometric method starts from such first principles and employs deductive reasoning to construct a comprehensive philosophical system based on them. Phenomenology gives particular importance to how things appear to be. It consists in suspending one's judgments about whether these things actually exist in the external world. This technique is known as epoché and can be used to study appearances independent of assumptions about their causes. The method of conceptual analysis came to particular prominence with the advent of analytic philosophy. It studies concepts by breaking them down into their most fundamental constituents to clarify their meaning. Common sense philosophy uses common and widely accepted beliefs as a philosophical tool. They are used to draw interesting conclusions. This is often employed in a negative sense to discredit radical philosophical positions that go against common sense. Ordinary language philosophy has a very similar method: it approaches philosophical questions by looking at how the corresponding terms are used in ordinary language.
Many methods in philosophy rely on some form of intuition. They are used, for example, to evaluate thought experiments, which involve imagining situations to assess their possible consequences in order to confirm or refute philosophical theories. The method of reflective equilibrium tries to form a coherent perspective by examining and reevaluating all the relevant beliefs and intuitions. Pragmatists focus on the practical consequences of philosophical theories to assess whether they are true or false. Experimental philosophy is a recently developed approach that uses the methodology of social psychology and the cognitive sciences for gathering empirical evidence and justifying philosophical claims.
Mathematics
In the field of mathematics, various methods can be distinguished, such as synthetic, analytic, deductive, inductive, and heuristic methods. For example, the difference between synthetic and analytic methods is that the former start from the known and proceed to the unknown while the latter seek to find a path from the unknown to the known. Geometry textbooks often proceed using the synthetic method. They start by listing known definitions and axioms and proceed by taking inferential steps, one at a time, until the solution to the initial problem is found. An important advantage of the synthetic method is its clear and short logical exposition. One disadvantage is that it is usually not obvious in the beginning that the steps taken lead to the intended conclusion. This may then come as a surprise to the reader since it is not explained how the mathematician knew in the beginning which steps to take. The analytic method often reflects better how mathematicians actually make their discoveries. For this reason, it is often seen as the better method for teaching mathematics. It starts with the intended conclusion and tries to find another formula from which it can be deduced. It then goes on to apply the same process to this new formula until it has traced back all the way to already proven theorems. The difference between the two methods concerns primarily how mathematicians think and present their proofs. The two are equivalent in the sense that the same proof may be presented either way.
Statistics
Statistics investigates the analysis, interpretation, and presentation of data. It plays a central role in many forms of quantitative research that have to deal with the data of many observations and measurements. In such cases, data analysis is used to cleanse, transform, and model the data to arrive at practically useful conclusions. There are numerous methods of data analysis. They are usually divided into descriptive statistics and inferential statistics. Descriptive statistics restricts itself to the data at hand. It tries to summarize the most salient features and present them in insightful ways. This can happen, for example, by visualizing its distribution or by calculating indices such as the mean or the standard deviation. Inferential statistics, on the other hand, uses this data based on a sample to draw inferences about the population at large. That can take the form of making generalizations and predictions or by assessing the probability of a concrete hypothesis.
Pedagogy
Pedagogy can be defined as the study or science of teaching methods. In this regard, it is the methodology of education: it investigates the methods and practices that can be applied to fulfill the aims of education. These aims include the transmission of knowledge as well as fostering skills and character traits. Its main focus is on teaching methods in the context of regular schools. But in its widest sense, it encompasses all forms of education, both inside and outside schools. In this wide sense, pedagogy is concerned with "any conscious activity by one person designed to enhance learning in another". The teaching happening this way is a process taking place between two parties: teachers and learners. Pedagogy investigates how the teacher can help the learner undergo experiences that promote their understanding of the subject matter in question.
Various influential pedagogical theories have been proposed. Mental-discipline theories were already common in ancient Greek and state that the main goal of teaching is to train intellectual capacities. They are usually based on a certain ideal of the capacities, attitudes, and values possessed by educated people. According to naturalistic theories, there is an inborn natural tendency in children to develop in a certain way. For them, pedagogy is about how to help this process happen by ensuring that the required external conditions are set up. Herbartianism identifies five essential components of teaching: preparation, presentation, association, generalization, and application. They correspond to different phases of the educational process: getting ready for it, showing new ideas, bringing these ideas in relation to known ideas, understanding the general principle behind their instances, and putting what one has learned into practice. Learning theories focus primarily on how learning takes place and formulate the proper methods of teaching based on these insights. One of them is apperception or association theory, which understands the mind primarily in terms of associations between ideas and experiences. On this view, the mind is initially a blank slate. Learning is a form of developing the mind by helping it establish the right associations. Behaviorism is a more externally oriented learning theory. It identifies learning with classical conditioning, in which the learner's behavior is shaped by presenting them with a stimulus with the goal of evoking and solidifying the desired response pattern to this stimulus.
The choice of which specific method is best to use depends on various factors, such as the subject matter and the learner's age. Interest and curiosity on the side of the student are among the key factors of learning success. This means that one important aspect of the chosen teaching method is to ensure that these motivational forces are maintained, through intrinsic or extrinsic motivation. Many forms of education also include regular assessment of the learner's progress, for example, in the form of tests. This helps to ensure that the teaching process is successful and to make adjustments to the chosen method if necessary.
Related concepts
Methodology has several related concepts, such as paradigm and algorithm. In the context of science, a paradigm is a conceptual worldview. It consists of a number of basic concepts and general theories, that determine how the studied phenomena are to be conceptualized and which scientific methods are considered reliable for studying them. Various theorists emphasize similar aspects of methodologies, for example, that they shape the general outlook on the studied phenomena and help the researcher see them in a new light.
In computer science, an algorithm is a procedure or methodology to reach the solution of a problem with a finite number of steps. Each step has to be precisely defined so it can be carried out in an unambiguous manner for each application. For example, the Euclidean algorithm is an algorithm that solves the problem of finding the greatest common divisor of two integers. It is based on simple steps like comparing the two numbers and subtracting one from the other.
| Physical sciences | Science basics | Basics and measurement |
1605738 | https://en.wikipedia.org/wiki/Tarsus%20%28skeleton%29 | Tarsus (skeleton) | In the human body, the tarsus (: tarsi) is a cluster of seven articulating bones in each foot situated between the lower end of the tibia and the fibula of the lower leg and the metatarsus. It is made up of the midfoot (cuboid, medial, intermediate, and lateral cuneiform, and navicular) and hindfoot (talus and calcaneus).
The tarsus articulates with the bones of the metatarsus, which in turn articulate with the proximal phalanges of the toes. The joint between the tibia and fibula above and the tarsus below is referred to as the ankle joint proper.
In humans the largest bone in the tarsus is the calcaneus, which is the weight-bearing bone within the heel of the foot.
Human anatomy
Bones
The talus bone or ankle bone is connected superiorly to the two bones of the lower leg, the tibia and fibula, to form the ankle joint or talocrural joint; inferiorly, at the subtalar joint, to the calcaneus or heel bone. Together, the talus and calcaneus form the hindfoot.
The five irregular bones of the midfoot—the cuboid, navicular, and three cuneiform bones—form the arches of the foot which serves as a shock absorber. The midfoot is connected to the hind- and forefoot by muscles and the plantar fascia.
Movements
The complex motion of the subtalar joint occurs in three planes and produces subtalar inversion and eversion. Along with the transverse tarsal joint (i.e. talonavicular and calcaneocuboid joint), the subtalar joint transforms tibial rotation into forefoot supination and pronation.
The axis of rotation in the joint is directed upward 42 degrees from the horizontal plane and 16 degrees medially from the midline of the foot.
However, together, the subtalar facets form a screw or Archimedean spiral (right-handed in the right foot) about which subtalar motion occurs. So, during subtalar inversion, the calcaneus also rotates clockwise and translates forward along the axis of the screw.
Average subtalar motion is 20-30 degrees inversion and 5-10 degrees eversion. Functional motion during the gait cycle is 10-15 degrees (the heel strikes the ground in slight inversion followed by quick eversion).
The talonavicular and calcaneocuboid joints (i.e. between the talus and navicular bones, and the calcaneus and cuboid bones) form the so-called transverse tarsal joint or Chopart's joint. It has two axes of motion.
Inversion and eversion occur about a longitudinal axis oriented 15 degrees upward from the horizontal plane and 9 degrees medially from the longitudinal axis of the foot.
Flexion and extension occur primarily about an oblique axis oriented 52 degrees upward from the horizontal plane and 57 degrees anteromedially (forward-inward).
In vitro talonavicular motion is 7 degrees flexion-extension and 17 degrees pronation-supination; while calcaneocuboid motion is 2 degrees flexion-extension and 7 degrees pronation-supination.
The motions of the subtalar and transverse talar joints interact to make the foot either flexible or rigid. With the subtalar joint in eversion, the two joints of the transverse joint are parallel, which make movements in this joint possible. With the subtalar joint in inversion, the axes of the transverse joint are convergent, movements in this joint are thus locked and the midfoot rigid.
Other animals
In primitive tetrapods, such as Trematops, the tarsus consists of three rows of bones. There are three proximal tarsals, the tibiale, intermedium, and fibulare, named for their points of articulation with the bones of the lower limb. These are followed by a second row of four bones, referred to as the centralia (singular: centrale), and then a row of five distal tarsals, each articulating with a single metatarsal. In the great majority of tetrapods, including all of those alive today, this simple pattern is modified by the loss and fusion of some of the bones.
In reptiles and mammals, there are normally just two proximal tarsals, the calcaneus (equivalent to the amphibian fibulare) and the talus (probably derived from a fusion of multiple bones). In mammals, including humans, the talus forms a hinge joint with the tibia, a feature especially well developed in the artiodactyls. The calcaneus is also modified, forming a heel for the attachment of the Achilles tendon. Neither of these adaptations is found in reptiles, which have a relatively simple structure to both bones.
The fifth distal tarsal disappears relatively early in evolution, with the remainder becoming the cuneiform and cuboid bones. Reptiles usually retain two centralia, while mammals typically have only one (the navicular).
In birds, the tarsus has disappeared, with the proximal tarsals having fused with the tibia, the centralia having disappeared, and the distal bones having fused with the metatarsals to form a single tarsometatarsus bone, effectively giving the leg a third segment.
Additional images
| Biology and health sciences | Skeletal system | Biology |
1605931 | https://en.wikipedia.org/wiki/Intermodal%20passenger%20transport | Intermodal passenger transport | Intermodal passenger transport, also called mixed-mode commuting, involves using two or more modes of transportation in a journey. Mixed-mode commuting is often used to combine the strengths (and offset the weaknesses) of various transportation options. A major goal of modern intermodal passenger transport is to reduce dependence on the automobile as the major mode of ground transportation and increase use of public transport. To assist the traveller, various intermodal journey planners such as Rome2rio and Google Transit have been devised to help travellers plan and schedule their journey.
Mixed-mode commuting often centers on one type of rapid transit, such as regional rail, to which low-speed options (i.e. bus, tram, or bicycle) are appended at the beginning or end of the journey. Trains offer quick transit from a suburb into an urban area, where passengers can choose a way to complete the trip. Most transportation modes have always been used intermodally; for example, people have used road or urban railway to an airport or inter-regional railway station.
History
Intermodal transport has existed for about as long as passenger transport itself. People switched from carriages to ferries at the edge of a river too deep to ford. In the 19th century, people who lived inland switched from train to ship for overseas voyages. Hoboken Terminal in Hoboken, New Jersey, was built to let commuters to New York City from New Jersey switch to ferries to cross the Hudson River in order to get to Manhattan. A massive ferry slip, now in ruins, was incorporated into the terminal building. Later, when a subway was built through tunnels under the Hudson, now called the PATH, a station stop was added to Hoboken Terminal. More recently, the New Jersey Transit's Hudson-Bergen Light Rail system has included a stop there. Ferry service has recently been revived, but passengers must exit the terminal and walk across the pier to the more modest ferry slip.
With the opening of the Woodside and Birkenhead Dock Street Tramway in 1873, Birkenhead Dock railway station probably became the world's first tram to train interchange station.
Urban mixed-mode commuting
Public transportation systems such as train or metro systems have the most efficient means and highest capacity to transport people around cities. Therefore, mixed-mode commuting in the urban environment is largely dedicated to first getting people onto the train network and once off the train network to their final destination.
Automobile to public transport nodes
Although automobiles are conventionally used as a single-mode form of transit, they also find use in a variety of mixed-mode scenarios. They can provide a short commute to train stations, airports, and piers, where all-day "park and ride" lots are often available. Used in this context, cars offer commuters the relative comfort of single-mode travel, while significantly reducing the financial and environmental costs.
Taxicabs and rental cars also play a major role in providing door-to-door service between airports or train stations and other points of travel throughout urban, suburban, and rural communities.
(Automobiles can also be used as the centerpiece of a multi-mode commute, with drivers resorting to walking or cycling to their final destination. Commuters to major cities take this route when driving is convenient, but parking options at the destination are not readily available.)
Park and ride
Transport planners often try to encourage automobile commuters to make much of their journey by public transport. One way of doing this is to provide car parking places at train or bus stations where commuters can drive to the station, park their cars and then continue on with their journey on the train or bus: this is often called "park and ride".
Similar to park and ride is what is often termed "kiss and ride". Rather than drive to the train or bus station and park the commuter is driven to the station by a friend or relative (parent, spouse etc.) The "kiss" refers to the peck on the cheek as the commuter exits the car. Kiss and ride is usually conducted when the train/bus/ferry station is close to home, so that the driver dropping the commuter off has a short journey to and from home.
Bus to public transport nodes
Many large cities link their railway network to their bus network. This enables commuters to get to places that are not serviced directly by rail as they are often considered to be too far for walking.
Feeder buses are a specific example of this; feeder buses service local neighbourhoods by taking travellers from their homes to nearby train stations which is important if the distances are too far to comfortably walk; at the end of the working day the buses take the travellers home again. Feeder buses work best when they are scheduled to arrive at the railway station shortly before the train arrives allowing enough time for commuters to comfortably walk to their train, and on the commuters' return journey buses are scheduled to arrive shortly after the train arrives so that the buses are waiting to take the commuters home. If train and bus services are very frequent then this scheduling is unimportant as the commuter will in any case have a very short wait to interchange.
Bike and ride
All around the world bicycles are used to get to and from train and other public transportation stations; this form of intermodal passenger transport is often called "bike and ride". To safeguard against theft or vandalism of parked bicycles at these train, bus, and ferry stations, "bike and ride" transport benefits greatly from secure bicycle parking facilities such as bicycle parking stations being available.
Some train, bus, and ferry systems allow commuters to take their bicycles aboard, allowing cyclists to ride at both ends of the commute, though sometimes this is restricted to off-peak travel periods: in such cases, folding bicycles may be permitted where regular bicycles are not. In some cities, bicycles are permitted aboard trains and buses.
In some cities a public bicycle rental programme allows commuters to take a public bike between the public transport station and a docking station near their origin or destination.
The use of "bike and ride" instead of a car can cut costs for fuel and parking, and some families no longer need to own and operate multiple cars. Environmental benefits can also increase (i.e. less pollution) and reduced traffic congestion can deliver significant cost savings to the city and local government.
Many transit agencies have begun installing bike racks on the front of buses, as well as in the interior of buses, trains, and even on ferries. These transit bike racks allow cyclists the ability to ride their bicycle to the bus/train/ferry, take the mode of transportation, then ride again to their final destination. These types of racks combined with increased bike infrastructure and bike parking have made bike commuting a frequent topic of discussion by cities and local government.
Inter-regional mixed mode commuting
Intermodal passenger transport involving air travel
Airport rail link
Many cities have extended subway or rail service to major urban airports. This provides travellers with an inexpensive, frequent and reliable way to get to their flights as opposed to driving or being driven, and contending with full up parking, or taking taxis and getting caught in traffic jams on the way to the airport. Many airports now have some mass transit link, including London, Sydney, Munich, Hong Kong, Vancouver, Philadelphia, Cleveland, New York City (JFK), Delhi, and Chennai.
Airport–ferry connection
At the Hong Kong International Airport, ferry services to various piers in the Pearl River Delta are provided. Passengers from Guangdong can use these piers to take a flight at the airport, without passing through customs and immigration control, effectively like having a transit from one flight to another. The airport is well-connected with expressways and an Airport Express train service. A seaport and logistics facilities will be added in the near future. Kansai International Airport is also connected to Kobe Airport with ferries. The Toronto Island ferry connects Billy Bishop Toronto City Airport to mainland Toronto, where passengers can connect to the Toronto streetcar system or with airport shuttle buses which transports to bus, subway and rail connections at Union Station.
Automobiles on trains
Several passenger rail systems offer services that allow travelers to bring their automobiles with them. These usually consist of automobile carrying wagons attached to normal passenger trains, but some special trains operate solely to transport automobiles. This is particularly of use in areas where trains may travel but automobiles cannot, such as the Channel Tunnel. Another system called NIMPR is designed to transport electric vehicles on high speed trains.
Trains on boats
A train ferry is a ship designed to carry railway vehicles. While usually used to carry freight vehicles, passenger cars can also be carried. In other places passengers move between passenger cars to a passenger ferry.
Train–ferry connection
Prior to the widespread use of automobiles, the San Francisco Bay Area featured a complex network of ferry services which connected numerous interurban and streetcar systems in the North and East Bay to the San Francisco Ferry Building, where several city streetcar lines began service. The opening of the rail-carrying San Francisco–Oakland Bay Bridge and automotive Golden Gate Bridge almost entirely supplanted these services.
Sonoma–Marin Area Rail Transit commuter rail is expected to feature a connection with the Golden Gate Ferry and service to San Francisco Ferry Building at Larkspur Landing. The Hercules station is to be the first direct Amtrak-to-ferry transit hub in the San Francisco Bay.
The Staten Island Railway, while operated by the Metropolitan Transportation Authority, does not have a physical connection to the rest of New York City's rail network. As such, transfers to Manhattan are facilitated by the free Staten Island Ferry.
Transfer facilities
In recent years, an increasing emphasis has been placed on designing facilities that make such transfers easier and more seamless. These are intended to help passengers move from one mode (or form) of transportation to another. An intermodal station may service air, rail, and highway transportation for example.
In some cases, facilities were merged or transferred into a new facility, as at the William F. Walsh Regional Transportation Center in Syracuse, New York, or South Station in Boston, Massachusetts. In other cases new facilities, such as the Alewife Station In Cambridge, Massachusetts, were built from the start to emphasize intermodalism.
Regional transit systems in the United States often include regional intermodal transit centers that incorporate multiple types of rail and bus services alongside park and ride amenities. Until the completion of San Francisco Salesforce Transit Center, the Millbrae Intermodal Terminal in California is the largest intermodal transit center west of the Mississippi which includes direct on-platform connections between BART, the Bay Area's regional rail system, Caltrain, the San Francisco Peninsula's commuter rail, and SamTrans, the regional bus service for San Mateo County. The uniqueness of this transfer facility is that turnstiles are located on the platforms between rail services in addition to on a separate concourse to allow for direct transfers. Millbrae Intermodal Terminal is also planned to be incorporated into the California High-Speed Rail project as one of two stations between San Francisco and San Jose.
Assessment
Advantages
Mixed mode commuting combines the benefits of walking, bicycle commuting, or driving with the benefits of rapid transit while offsetting some of the major disadvantages of each. The use of a bicycle can, for example, make an (inexpensive compared to a car) 20 mile light-rail or suburban rail journey attractive even if the endpoints of the journey each sit 1 mile out from the stations: the 30 minutes walking time becomes 8 minutes bicycling.
As in the example above, location plays a large role in mixed mode commuting. Rapid transit such as express bus or light rail may cover most of the distance, but sit too far out from commute endpoints. At 3 mph walking, 2 miles represents about 40 minutes of commute time; whereas a bicycle may pace 12 mph leisurely, cutting this time to 10 minutes. When the commuter finds the distance between the originating endpoint (e.g. the home) and the destination (e.g. the place of employment) too far to be enjoyable or practical, commute by car or motorcycle to the station may remain practical, as long as the commute from the far end station to the destination is practical by walking, a carry-on cycle, or another rapid transit such as a local or shuttle bus.
In general, locations close to major transit such as rail stations carry higher land value and thus higher costs to rent or purchase. A commuter may select a location further out than practical walking distance but not more than practical cycling distance to reduce housing costs. Similarly, a commuter can close an even further distance quickly with an ebike, motorcycle, or car, allowing for the selection of a more preferred living area somewhat further from the station than would be viable by walking or simple bicycle.
Other cost advantages of mixed mode commuting include lower vehicle insurance via Pay As You Drive programs; lower fuel and maintenance costs; and increased automobile life. In the most extreme cases, a mixed-mode commuter may opt to car share and pay only a small portion of purchase, fuel, maintenance, and insurance, or to live car-free. These cost benefits are offset by costs of transit, which can vary. A Maryland MTA month pass valid for MTA Light Rail, Metro Subway, and City Bus costs $64, while a month pass for the Baltimore to DC MARC costs $175.00 and a DC MetroRail 7 day pass costs $47 totaling $182. In most of Europe :de:Verkehrsverbund and mode neutral pricing eliminate the need to have several different tickets for public transit across different modes. Mobility as a service intends to take this a step further, offering one price per trip from door to door, no matter which mode is used for which part of the trip.
Disadvantages
The effectiveness of a mixed-mode commute can be measured in many ways: speed to destination, convenience, security, environmental impact, and proximity to mass transit are all factors. Because mixed-mode commutes rely on a certain degree of coordination, scheduling issues with mass transit can often be an issue. For example, a sometimes-late train can be an annoyance, and an often-late train can make a commute impractical.
Weather can also be a factor. Even when the use of an automobile is involved, the transition from one mode of transportation to another often exposes commuters to the elements. As a result, multi-mode commuters often travel prepared for inclement weather.
In the United States fare integration is often lacking, making passengers "pay extra for the 'privilege' of having a connection". This is largely a non-issue in European cities where all modes of local public transit follow the same ticketing scheme and a ticket for e.g. the metro will be valid on buses or commuter rail.
| Technology | Basics_7 | null |
1607203 | https://en.wikipedia.org/wiki/FM%20broadcasting | FM broadcasting | FM broadcasting is a method of radio broadcasting that uses frequency modulation (FM) of the radio broadcast carrier wave. Invented in 1933 by American engineer Edwin Armstrong, wide-band FM is used worldwide to transmit high-fidelity sound over broadcast radio. FM broadcasting offers higher fidelity—more accurate reproduction of the original program sound—than other broadcasting techniques, such as AM broadcasting. It is also less susceptible to common forms of interference, having less static and popping sounds than are often heard on AM. Therefore, FM is used for most broadcasts of music and general audio (in the audio spectrum). FM radio stations use the very high frequency range of radio frequencies.
Broadcast bands
Throughout the world, the FM broadcast band falls within the VHF part of the radio spectrum. Usually 87.5 to 108.0 MHz is used, or some portion of it, with few exceptions:
In the former Soviet republics, and some former Eastern Bloc countries, the older 65.8–74 MHz band is also used. Assigned frequencies are at intervals of 30 kHz. This band, sometimes referred to as the OIRT band, is slowly being phased out. Where the OIRT band is used, the 87.5–108.0 MHz band is referred to as the CCIR band.
In Japan, the band 76–95 MHz is used.
In Brazil, until the late 2010s, FM broadcast stations only used the 88–108 MHz band, but with the phasing out of analog television, the 76-88 MHz band (old band channels 5 and 6 in VHF television) are allocated for old local MW stations which have moved to FM in agreement with ANATEL.
The frequency of an FM broadcast station (more strictly its assigned nominal center frequency) is usually a multiple of 100 kHz. In most of South Korea, the Americas, the Philippines, and the Caribbean, only odd multiples are used. Some other countries follow this plan because of the import of vehicles, principally from the United States, with radios that can only tune to these frequencies. In some parts of Europe, Greenland, and Africa, only even multiples are used. In the United Kingdom, both odd and even are used. In Italy, multiples of 50 kHz are used. In most countries the maximum permitted frequency error of the unmodulated carrier is specified, which typically should be within 2 kHz of the assigned frequency.
There are other unusual and obsolete FM broadcasting standards in some countries, with non-standard spacings of 1, 10, 30, 74, 500, and 300 kHz. To minimise inter-channel interference, stations operating from the same or nearby transmitter sites tend to keep to at least a 500 kHz frequency separation even when closer frequency spacing is technically permitted. The ITU publishes Protection Ratio graphs, which give the minimum spacing between frequencies based on their relative strengths. Only broadcast stations with large enough geographic separations between their coverage areas can operate on the same or close frequencies.
Technology
Modulation
Frequency modulation or FM is a form of modulation which conveys information by varying the frequency of a carrier wave; the older amplitude modulation or AM varies the amplitude of the carrier, with its frequency remaining constant. With FM, frequency deviation from the assigned carrier frequency at any instant is directly proportional to the amplitude of the (audio) input signal, determining the instantaneous frequency of the transmitted signal. Because transmitted FM signals use significantly more bandwidth than AM signals, this form of modulation is commonly used with the higher (VHF or UHF) frequencies used by TV, the FM broadcast band, and land mobile radio systems.
The maximum frequency deviation of the carrier is usually specified and regulated by the licensing authorities in each country. For a stereo broadcast, the maximum permitted carrier deviation is invariably ±75 kHz, although a little higher is permitted in the United States when SCA systems are used. For a monophonic broadcast, again the most common permitted maximum deviation is ±75 kHz. However, some countries specify a lower value for monophonic broadcasts, such as ±50 kHz.
Bandwidth
The bandwidth of an FM transmission is given by the Carson bandwidth rule which is the sum of twice the maximum deviation and twice the maximum modulating frequency. For a transmission that includes RDS this would be = . This is also known as the necessary bandwidth.
Pre-emphasis and de-emphasis
Random noise has a triangular spectral distribution in an FM system, with the effect that noise occurs predominantly at the higher audio frequencies within the baseband. This can be offset, to a limited extent, by boosting the high frequencies before transmission and reducing them by a corresponding amount in the receiver. Reducing the high audio frequencies in the receiver also reduces the high-frequency noise. These processes of boosting and then reducing certain frequencies are known as pre-emphasis and de-emphasis, respectively.
The amount of pre-emphasis and de-emphasis used is defined by the time constant of a simple RC filter circuit. In most of the world a 50 μs time constant is used. In the Americas and South Korea, 75 μs is used. This applies to both mono and stereo transmissions. For stereo, pre-emphasis is applied to the left and right channels before multiplexing.
The use of pre-emphasis becomes a problem because many forms of contemporary music contain more high-frequency energy than the musical styles which prevailed at the birth of FM broadcasting. Pre-emphasizing these high-frequency sounds would cause excessive deviation of the FM carrier. Modulation control (limiter) devices are used to prevent this. Systems more modern than FM broadcasting tend to use either programme-dependent variable pre-emphasis; e.g., dbx in the BTSC TV sound system, or none at all.
Pre-emphasis and de-emphasis was used in the earliest days of FM broadcasting. According to a BBC report from 1946, 100 μs was originally considered in the US, but 75 μs subsequently adopted.
Stereo FM
Long before FM stereo transmission was considered, FM multiplexing of other types of audio-level information was experimented with. Edwin Armstrong, who invented FM, was the first to experiment with multiplexing, at his experimental 41 MHz station W2XDG located on the 85th floor of the Empire State Building in New York City.
These FM multiplex transmissions started in November 1934 and consisted of the main channel audio program and three subcarriers: a fax program, a synchronizing signal for the fax program and a telegraph order channel. These original FM multiplex subcarriers were amplitude modulated.
Two musical programs, consisting of both the Red and Blue Network program feeds of the NBC Radio Network, were simultaneously transmitted using the same system of subcarrier modulation as part of a studio-to-transmitter link system. In April 1935, the AM subcarriers were replaced by FM subcarriers, with much improved results.
The first FM subcarrier transmissions emanating from Major Armstrong's experimental station KE2XCC at Alpine, New Jersey occurred in 1948. These transmissions consisted of two-channel audio programs, binaural audio programs and a fax program. The original subcarrier frequency used at KE2XCC was 27.5 kHz. The IF bandwidth was ±5 kHz, as the only goal at the time was to relay AM radio-quality audio. This transmission system used 75 μs audio pre-emphasis like the main monaural audio and subsequently the multiplexed stereo audio.
In the late 1950s, several systems to add stereo to FM radio were considered by the FCC. Included were systems from 14 proponents including Crosby, Halstead, Electrical and Musical Industries, Ltd (EMI), Zenith, and General Electric. The individual systems were evaluated for their strengths and weaknesses during field tests in Uniontown, Pennsylvania, using KDKA-FM in Pittsburgh as the originating station. The Crosby system was rejected by the FCC because it was incompatible with existing subsidiary communications authorization (SCA) services which used various subcarrier frequencies including 41 and 67 kHz. Many revenue-starved FM stations used SCAs for "storecasting" and other non-broadcast purposes. The Halstead system was rejected due to lack of high frequency stereo separation and reduction in the main channel signal-to-noise ratio. The GE and Zenith systems, so similar that they were considered theoretically identical, were formally approved by the FCC in April 1961 as the standard stereo FM broadcasting method in the United States and later adopted by most other countries. It is important that stereo broadcasts be compatible with mono receivers. For this reason, the left (L) and right (R) channels are algebraically encoded into sum (L+R) and difference (L−R) signals. A mono receiver will use just the L+R signal so the listener will hear both channels through the single loudspeaker. A stereo receiver will add the difference signal to the sum signal to recover the left channel, and subtract the difference signal from the sum to recover the right channel.
The (L+R) signal is limited to 30 Hz to 15 kHz to protect a 19 kHz pilot signal. The (L−R) signal, which is also limited to 15 kHz, is amplitude modulated onto a 38 kHz double-sideband suppressed-carrier (DSB-SC) signal, thus occupying 23 kHz to 53 kHz. A 19 kHz ± 2 Hz pilot tone, at exactly half the 38 kHz sub-carrier frequency and with a precise phase relationship to it, as defined by the formula below, is also generated. The pilot is transmitted at 8–10% of overall modulation level and used by the receiver to identify a stereo transmission and to regenerate the 38 kHz sub-carrier with the correct phase. The composite stereo multiplex signal contains the Main Channel (L+R), the pilot tone, and the (L−R) difference signal. This composite signal, along with any other sub-carriers, modulates the FM transmitter. The terms composite, multiplex and even MPX are used interchangeably to describe this signal.
The instantaneous deviation of the transmitter carrier frequency due to the stereo audio and pilot tone (at 10% modulation) is
where A and B are the pre-emphasized left and right audio signals and =19 kHz is the frequency of the pilot tone. Slight variations in the peak deviation may occur in the presence of other subcarriers or because of local regulations.
Another way to look at the resulting signal is that it alternates between left and right at 38 kHz, with the phase determined by the 19 kHz pilot signal. Most stereo encoders use this switching technique to generate the 38 kHz subcarrier, but practical encoder designs need to incorporate circuitry to deal with the switching harmonics. Converting the multiplex signal back into left and right audio signals is performed by a decoder, built into stereo receivers. Again, the decoder can use a switching technique to recover the left and right channels.
In addition, for a given RF level at the receiver, the signal-to-noise ratio and multipath distortion for the stereo signal will be worse than for the mono receiver. For this reason many stereo FM receivers include a stereo/mono switch to allow listening in mono when reception conditions are less than ideal, and most car radios are arranged to reduce the separation as the signal-to-noise ratio worsens, eventually going to mono while still indicating a stereo signal is received. As with monaural transmission, it is normal practice to apply pre-emphasis to the left and right channels before encoding and to apply de-emphasis at the receiver after decoding.
In the U.S. around 2010, using single-sideband modulation for the stereo subcarrier was proposed. It was theorized to be more spectrum-efficient and to produce a 4 dB s/n improvement at the receiver, and it was claimed that multipath distortion would be reduced as well. A handful of radio stations around the country broadcast stereo in this way, under FCC experimental authority. It may not be compatible with very old receivers, but it is claimed that no difference can be heard with most newer receivers. At present, the FCC rules do not allow this mode of stereo operation.
Quadraphonic FM
In 1969, Louis Dorren invented the Quadraplex system of single station, discrete, compatible four-channel FM broadcasting. There are two additional subcarriers in the Quadraplex system, supplementing the single one used in standard stereo FM. The baseband layout is as follows:
50 Hz to 15 kHz main channel (sum of all 4 channels) (LF+LR+RF+RR) signal, for mono FM listening compatibility.
23 to 53 kHz (sine quadrature subcarrier) (LF+LR) − (RF+RR) left minus right difference signal. This signal's modulation in algebraic sum and difference with the main channel is used for 2 channel stereo listener compatibility.
23 to 53 kHz (cosine quadrature 38 kHz subcarrier) (LF+RR) − (LR+RF) Diagonal difference. This signal's modulation in algebraic sum and difference with the main channel and all the other subcarriers is used for the Quadraphonic listener.
61 to 91 kHz (sine quadrature 76 kHz subcarrier) (LF+RF) − (LR+RR) Front-back difference. This signal's modulation in algebraic sum and difference with the main channel and all the other subcarriers is also used for the Quadraphonic listener.
105 kHz SCA subcarrier, phase-locked to 19 kHz pilot, for reading services for the blind, background music, etc.
The normal stereo signal can be considered as switching between left and right channels at 38 kHz, appropriately band-limited. The quadraphonic signal can be considered as cycling through LF, LR, RF, RR, at 76 kHz.
Early efforts to transmit discrete four-channel quadraphonic music required the use of two FM stations; one transmitting the front audio channels, the other the rear channels. A breakthrough came in 1970 when KIOI (K-101) in San Francisco successfully transmitted true quadraphonic sound from a single FM station using the Quadraplex system under Special Temporary Authority from the FCC. Following this experiment, a long-term test period was proposed that would permit one FM station in each of the top 25 U.S. radio markets to transmit in Quadraplex. The test results hopefully would prove to the FCC that the system was compatible with existing two-channel stereo transmission and reception and that it did not interfere with adjacent stations.
There were several variations on this system submitted by GE, Zenith, RCA, and Denon for testing and consideration during the National Quadraphonic Radio Committee field trials for the FCC. The original Dorren Quadraplex System outperformed all the others and was chosen as the national standard for Quadraphonic FM broadcasting in the United States. The first commercial FM station to broadcast quadraphonic program content was WIQB (now called WWWW-FM) in Ann Arbor/Saline, Michigan under the guidance of Chief Engineer Brian Jeffrey Brown.
Noise reduction
Various attempts to add analog noise reduction to FM broadcasting were carried out in the 1970s and 1980s:
A commercially unsuccessful noise reduction system used with FM radio in some countries during the late 1970s, Dolby FM was similar to Dolby B but used a modified 25 μs pre-emphasis time constant and a frequency selective companding arrangement to reduce noise. The pre-emphasis change compensates for the excess treble response that otherwise would make listening difficult for those without Dolby decoders.
A similar system named High Com FM was tested in Germany between July 1979 and December 1981 by IRT. It was based on the Telefunken High Com broadband compander system, but was never introduced commercially in FM broadcasting.
Yet another system was the CX-based noise reduction system FMX implemented in some radio broadcasting stations in the United States in the 1980s.
Other subcarrier services
FM broadcasting has included subsidiary communications authorization (SCA) services capability since its inception, as it was seen as another service which licensees could use to create additional income. Use of SCAs was particularly popular in the US, but much less so elsewhere. Uses for such subcarriers include radio reading services for the blind, which became common and remain so, private data transmission services (for example sending stock market information to stockbrokers or stolen credit card number denial lists to stores,) subscription commercial-free background music services for shops, paging ("beeper") services, alternative-language programming, and providing a program feed for AM transmitters of AM/FM stations. SCA subcarriers are typically 67 kHz and 92 kHz. Initially the users of SCA services were private analog audio channels which could be used internally or leased, for example Muzak-type services. There were experiments with quadraphonic sound. If a station does not broadcast in stereo, everything from 23 kHz on up can be used for other services. The guard band around 19 kHz (±4 kHz) must still be maintained, so as not to trigger stereo decoders on receivers. If there is stereo, there will typically be a guard band between the upper limit of the DSBSC stereo signal (53 kHz) and the lower limit of any other subcarrier.
Digital data services are also available. A 57 kHz subcarrier (phase locked to the third harmonic of the stereo pilot tone) is used to carry a low-bandwidth digital Radio Data System signal, providing extra features such as station name, alternative frequency (AF), traffic data for satellite navigation systems and radio text (RT). This narrowband signal runs at only 1,187.5 bits per second, thus is only suitable for text. A few proprietary systems are used for private communications. A variant of RDS is the North American RBDS. In Germany the analog ARI system was used prior to RDS to alert motorists that traffic announcements were broadcast (without disturbing other listeners). Plans to use ARI for other European countries led to the development of RDS as a more powerful system. RDS is designed to be capable of use alongside ARI despite using identical subcarrier frequencies.
In the United States and Canada, digital radio services are deployed within the FM band rather than using Eureka 147 or the Japanese standard ISDB. This in-band on-channel approach, as do all digital radio techniques, makes use of advanced compressed audio. The proprietary iBiquity system, branded as HD Radio, is authorized for "hybrid" mode operation, wherein both the conventional analog FM carrier and digital sideband subcarriers are transmitted.
Transmission power
The output power of an FM broadcasting transmitter is one of the parameters that governs how far a transmission will cover. The other important parameters are the height of the transmitting antenna and the antenna gain. Transmitter powers should be carefully chosen so that the required area is covered without causing interference to other stations further away. Practical transmitter powers range from a few milliwatts to 80 kW. As transmitter powers increase above a few kilowatts, the operating costs become high and only viable for large stations. The efficiency of larger transmitters is now better than 70% (AC power in to RF power out) for FM-only transmission. This compares to 50% before high efficiency switch-mode power supplies and LDMOS amplifiers were used. Efficiency drops dramatically if any digital HD Radio service is added.
Reception distance
VHF radio waves usually do not travel far beyond the visual horizon, so reception distances for FM stations are typically limited to . They can also be blocked by hills and to a lesser extent by buildings. Individuals with more-sensitive receivers or specialized antenna systems, or who are located in areas with more favorable topography, may be able to receive useful FM broadcast signals at considerably greater distances.
The knife edge effect can permit reception where there is no direct line of sight between broadcaster and receiver. The reception can vary considerably depending on the position. One example is the Učka mountain range, which makes constant reception of Italian signals from Veneto and Marche possible in a good portion of Rijeka, Croatia, despite the distance being over 200 km (125 miles). Other radio propagation effects such as tropospheric ducting and Sporadic E can occasionally allow distant stations to be intermittently received over very large distances (hundreds of miles), but cannot be relied on for commercial broadcast purposes. Good reception across the country is one of the main advantages over DAB/+ radio.
This is still less than the range of AM radio waves, which because of their lower frequencies can travel as ground waves or reflect off the ionosphere, so AM radio stations can be received at hundreds (sometimes thousands) of miles. This is a property of the carrier wave's typical frequency (and power), not its mode of modulation.
The range of FM transmission is related to the transmitter's RF power, the antenna gain, and antenna height. Interference from other stations is also a factor in some places. In the U.S, the FCC publishes curves that aid in calculation of this maximum distance as a function of signal strength at the receiving location. Computer modelling is more commonly used for this around the world.
Many FM stations, especially those located in severe multipath areas, use extra audio compression/processing to keep essential sound above the background noise for listeners, often at the expense of overall perceived sound quality. In such instances, however, this technique is often surprisingly effective in increasing the station's useful range.
History
Americas
Brazil
The first radio station to broadcast in FM in Brazil was Rádio Imprensa, which began broadcasting in Rio de Janeiro in 1955, on the 102.1 MHz frequency, founded by businesswoman Anna Khoury. Due to the high import costs of FM radio receivers, transmissions were carried out in circuit closed to businesses and stores, which played ambient music offered by radio. Until 1976, Rádio Imprensa was the only station operating in FM in Brazil. From the second half of the 1970s onwards, FM radio stations began to become popular in Brazil, causing AM radio to gradually lose popularity.
In 2021, the Brazilian Ministry of Communications expanded the FM radio band from 87.5-108.0 MHz to 76.1-108.0 MHz to enable the migration of AM radio stations in Brazilian capitals and large cities.
United States
FM broadcasting began in the late 1930s, when it was initiated by a handful of early pioneer experimental stations, including W1XOJ/W43B/WGTR (shut down in 1953) and W1XTG/WSRS, both transmitting from Paxton, Massachusetts (now listed as Worcester, Massachusetts); W1XSL/W1XPW/W65H/WDRC-FM/WFMQ/WHCN, Meriden, Connecticut; and W2XMN, KE2XCC, and WFMN, Alpine, New Jersey (owned by Edwin Armstrong himself, closed down upon Armstrong's death in 1954). Also of note were General Electric stations W2XDA Schenectady and W2XOY New Scotland, New York—two experimental FM transmitters on 48.5 MHz—which signed on in 1939. The two began regular programming, as W2XOY, on November 20, 1940. Over the next few years this station operated under the call signs W57A, W87A and WGFM, and moved to 99.5 MHz when the FM band was relocated to the 88–108 MHz portion of the radio spectrum. General Electric sold the station in the 1980s. Today this station is WRVE.
Other pioneers included W2XQR/W59NY/WQXQ/WQXR-FM, New York; W47NV/WSM-FM Nashville, Tennessee (signed off in 1951); W1XER/W39B/WMNE, with studios in Boston and later Portland, Maine, but whose transmitter was atop the highest mountain in the northeast United States, Mount Washington, New Hampshire (shut down in 1948); and W9XAO/W55M/WTMJ-FM Milwaukee, Wisconsin (went off air in 1950).
A commercial FM broadcasting band was formally established in the United States as of January 1, 1941, with the first fifteen construction permits announced on October 31, 1940. These stations primarily simulcast their AM sister stations, in addition to broadcasting lush orchestral music for stores and offices, classical music to an upmarket listenership in urban areas, and educational programming.
On June 27, 1945 the FCC announced the reassignment of the FM band to 90 channels from 88–106 MHz (which was soon expanded to 100 channels from 88–108 MHz). This shift, which the AM-broadcaster RCA had pushed for, made all the Armstrong-era FM receivers useless and delayed the expansion of FM. In 1961 WEFM (in the Chicago area) and WGFM (in Schenectady, New York) were reported as the first stereo stations. By the late 1960s, FM had been adopted for broadcast of stereo "A.O.R.—'Album Oriented Rock' Format", but it was not until 1978 that listenership to FM stations exceeded that of AM stations in North America. In most of the 70s FM was seen as highbrow radio associated with educational programming and classical music, which changed during the 1980s and 1990s when Top 40 music stations and later even country music stations largely abandoned AM for FM. Today AM is mainly the preserve of talk radio, news, sports, religious programming, ethnic (minority language) broadcasting and some types of minority interest music. This shift has transformed AM into the "alternative band" that FM once was. (Some AM stations have begun to simulcast on, or switch to, FM signals to attract younger listeners and aid reception problems in buildings, during thunderstorms, and near high-voltage wires. Some of these stations now emphasize their presence on the FM band.)
Europe
The medium wave band (known as the AM band because most stations using it employ amplitude modulation) was overcrowded in western Europe, leading to interference problems and, as a result, many MW frequencies are suitable only for speech broadcasting.
Belgium, the Netherlands, Denmark and particularly Germany were among the first countries to adopt FM on a widespread scale. Among the reasons for this were:
The medium wave band in Western Europe became overcrowded after World War II, mainly due to the best available medium wave frequencies used at high power levels by the Allied Occupation Forces, both for broadcasting entertainment to their troops and for broadcasting Cold War propaganda across the Iron Curtain.
After World War II, broadcasting frequencies were reorganized and reallocated by delegates of the victorious countries in the Copenhagen Frequency Plan. German broadcasters were left with only two remaining AM frequencies and were forced to look to FM for expansion.
Public service broadcasters in Ireland and Australia were far slower at adopting FM radio than those in either North America or continental Europe.
Netherlands
Hans Idzerda operated a broadcasting station, PCGG, at The Hague from 1919 to 1924, which employed narrow-band FM transmissions.
United Kingdom
In the United Kingdom the BBC conducted tests during the 1940s, then began FM broadcasting in 1955, with three national networks: the Light Programme, Third Programme and Home Service. These three networks used the sub-band 88.0–94.6 MHz. The sub-band 94.6–97.6 MHz was later used for BBC and local commercial services. Experimental stereo broadcasts started in the London area in January 1958.
However, only when commercial broadcasting was introduced to the UK in 1973 did the use of FM pick up in Britain. With the gradual clearance of other users (notably Public Services such as police, fire and ambulance) and the extension of the FM band to 108.0 MHz between 1980 and 1995, FM expanded rapidly throughout the British Isles and effectively took over from LW and MW as the delivery platform of choice for fixed and portable domestic and vehicle-based receivers. In addition, Ofcom (previously the Radio Authority) in the UK issues on demand Restricted Service Licences on FM and also on AM (MW) for short-term local-coverage broadcasting which is open to anyone who does not carry a prohibition and can put up the appropriate licensing and royalty fees. In 2010 around 450 such licences were issued.
When the BBC's radio networks were renamed Radio 2, Radio 3 and Radio 4 respectively in 1967 to coincide with the launch of Radio 1, the new station was the only one of the main four to not have an FM frequency allocated, which was the case for 21 years. Instead, Radio 1 shared airtime with Radio 2 FM, on Saturday afternoons, Sunday evenings, weekday evenings (10pm to midnight) and Bank Holidays, eventually having its own FM frequency starting in London in October 1987 on 104.8 MHz from Crystal Palace. Eventually in 1987, a frequency range of 97.6-99.8 MHz was allocated once police mobile radio transmitters were moved from band II, starting in London before being nationally completed by 1989. Radio 1 in London moved from its previous frequency to 98.8 MHz transmitted from the BBC's Wrotham site in Kent. Following this the BBC Radio 1 FM frequencies were rolled out to the rest of the UK.
Italy
Italy adopted FM broadcast widely in the early 1970s, but first experiments made by RAI dated back to 1950, when the "movement for free radio", developed by so-called "pirates", forced the recognition of free speech rights also through the use of "free radio media such as Broadcast transmitters", and took the case to the Constitutional Court of Italy. The court finally decided in favor of Free Radio. Just weeks after the court's final decision there was an "FM radio boom" involving small private radio stations across the country. By the mid-1970s, every city in Italy had a crowded FM radio spectrum.
Greece
Greece was another European country where the FM radio spectrum was used at first by the so-called "pirates" (both in Athens and Thessaloniki, the two major Greek cities) in the mid-1970s, before any national stations had started broadcasting on it; there were many AM (MW) stations in use for the purpose. No later than the end of 1977, the national public service broadcasting company EIRT (later also known as ERT) placed in service its first FM transmitter in the capital, Athens. By the end of the 1970s, most of Greek territory was covered by three National FM programs, and every city had many FM "pirates" as well.
The adaptation of the FM band for privately owned commercial radio stations came far later, in 1987.
Australia
FM broadcasting started in Australian capital cities in 1947 on an "experimental" basis, using an ABC national network feed, consisting largely of classical music and Parliament, as a programme source. It had a very small audience and was shut down in 1961 ostensibly to clear the television band: TV channel 5 (102.250 video carrier) if allocated would fall within the VHF FM band (98–108 MHz). The official policy on FM at the time was to eventually introduce it on another band, which would have required FM tuners custom-built for Australia. This policy was finally reversed and FM broadcasting was reopened in 1975 using the VHF band, after the few encroaching TV stations had been moved. Subsequently, it developed steadily until in the 1980s many AM stations transferred to FM due to its superior sound quality and lower operating costs. Today, as elsewhere in the developed world, most urban Australian broadcasting is on FM, although AM talk stations are still very popular. Regional broadcasters still commonly operate AM stations due to the additional range the broadcasting method offers. Some stations in major regional centres simulcast on AM and FM bands. Digital radio using the DAB+ standard has been rolled out to capital cities.
New Zealand
Like Australia, New Zealand adopted the FM format relatively late. As was the case with privately owned AM radio in the late 1960s, it took a spate of 'pirate' broadcasters to persuade a control-oriented, technology-averse government to allow FM to be introduced after at least five years of consumer campaigning starting in the mid-1970s, particularly in Auckland. An experimental FM station, FM 90.7, was broadcast in Whakatāne in early 1982. Later that year, Victoria University of Wellington's Radio Active began full-time FM transmissions. Commercial FM licences were finally approved in 1983, with Auckland-based 91FM and 89FM being the first to take up the offer. Broadcasting was deregulated in 1989.
Like many other countries in Africa and Asia that drive on the left, New Zealand imports vehicles from Japan. The standard radios in these vehicles operate on 76-to-90 MHz, which is not compatible with the 88-to-108 MHz range. Imported cars with Japanese radios can have FM expanders installed which down-convert the higher frequencies above 90 MHz. New Zealand has no indigenous car manufacturers.
Trinidad and Tobago
Trinidad and Tobago's first FM Radio station was 95.1FM, now rebranded as 951 Remix, which was launched in March 1976 by the TBC Radio Network.
Turkey
In Turkey, FM broadcasting began in the late 1960s, carrying several shows from the One television network which was transferred from the AM frequency (also known as MW in Turkey). In subsequent years, more MW stations were slowly transferred to FM, and by the end of the 1970s, most radio stations that were previously on MW had been moved to FM, though many talk, news and sport, but mostly religious stations, still remain on MW.
Other countries
Most other countries implemented FM broadcasting through the 1960s and expanded their use of FM through the 1990s. Because it takes a large number of FM transmitting stations to cover a geographically large country, particularly where there are terrain difficulties, FM is more suited to local broadcasting than for national networks. In such countries, particularly where there are economic or infrastructural problems, "rolling out" a national FM broadcast network to reach the majority of the population can be a slow and expensive process. Despite this, mostly in east European countries, national FM broadcast networks were established in the late 1960s and 1970s. In all Soviet-dependent countries except GDR, the OIRT band was used. First restricted to 68–73 MHz with 100 kHz channel spacing, then in the 1970s eventually expanded to 65.84–74.00 MHz with 30 kHz channel spacing.
The use of FM for domestic radio encouraged listeners to acquire cheap FM-only receivers and so reduced the number able to listen to longer-range AM foreign broadcasters. Similar considerations led to domestic radio in South Africa switching to FM in the 1960s.
ITU Conferences about FM
The frequencies available for FM were decided by some important conferences of ITU. The milestone of those conferences is the Stockholm agreement of 1961 among 38 countries. A 1984 conference in Geneva made some modifications to the original Stockholm agreement particularly in the frequency range above 100 MHz.
FM broadcasting switch-off
In 2017, Norway became the first country to completely switch to Digital audio broadcasting, the exception being some local stations remaining on FM until 2022, and might be extended to 2031. The switchover to DAB+ meant that especially rural areas obtained a far more diverse radio content compared to the FM-only period; several new radio stations had started transmissions on DAB+ in the years before the FM switch-off.
Switzerland is in the process of becoming the second country to switch from FM to DAB+. Public broadcaster SRG SSR shut down its entire FM infrastructure on 31 December 2024, citing low usage of FM (estimated to be 10% of the audience and falling) from widespread adoption of DAB+ and the cost of maintaining two broadcast infrastructures in parallel. Private broadcasters are undertaking a gradual shutdown of their FM transmitters to be completed by 31 December 2026.
Small-scale use of the FM broadcast band
Consumer use of FM transmitters
In some countries, small-scale (Part 15 in United States terms) transmitters are available that can transmit a signal from an audio device (usually an MP3 player or similar) to a standard FM radio receiver; such devices range from small units built to carry audio to a car radio with no audio-in capability (often formerly provided by special adapters for audio cassette decks, which are no longer common on car radio designs) up to full-sized, near-professional-grade broadcasting systems that can be used to transmit audio throughout a property, including systems that synchronize holiday decorative lighting with music. Most such units transmit in full stereo, though some models designed for beginner hobbyists might not. Similar transmitters are often included in satellite radio receivers and some toys.
Legality of these devices varies by country. The U.S. Federal Communications Commission and Industry Canada allow them. Starting on 1 October 2006, these devices became legal in most countries in the European Union. Devices made to the harmonized European specification became legal in the UK on 8 December 2006.
The FM broadcast band is also used by some inexpensive wireless microphones sold as toys for karaoke or similar purposes, allowing the user to use an FM radio as an output rather than a dedicated amplifier and speaker. Professional-grade wireless microphones generally use bands in the UHF region so they can run on dedicated equipment without broadcast interference.
Some wireless headphones transmit in the FM broadcast band, with the headphones tunable to only a subset of the broadcast band. Higher-quality wireless headphones use infrared transmission or UHF ISM bands such as 315 MHz, 863 MHz, 915 MHz, or 2.4 GHz instead of the FM broadcast band.
Assistive listening
Some assistive listening devices are based on FM radio, mostly using the 72.1 to 75.8 MHz band. Aside from the assisted listening receivers, only certain kinds of FM receivers can tune to this band.
Microbroadcasting
Low-power transmitters such as those mentioned above are also sometimes used for neighborhood or campus radio stations, though campus radio stations are often run over carrier current. This is generally considered a form of microbroadcasting. As a general rule, enforcement towards low-power FM stations is stricter than with AM stations, due to problems such as the capture effect, and as a result, FM microbroadcasters generally do not reach as far as their AM competitors.
Clandestine use of FM transmitters
FM transmitters have been used to construct miniature wireless microphones for espionage and surveillance purposes (covert listening devices or so-called "bugs"); the advantage to using the FM broadcast band for such operations is that the receiving equipment would not be considered particularly suspect. Common practice is to tune the bug's transmitter off the ends of the broadcast band, into what in the United States would be TV channel 6 (<87.9 MHz) or aviation navigation frequencies (>107.9 MHz); most FM radios with analog tuners have sufficient overcoverage to pick up these slightly-beyond-outermost frequencies, although many digitally tuned radios have not.
Constructing a "bug" is a common early project for electronics hobbyists, and project kits to do so are available from a wide variety of sources. The devices constructed, however, are often too large and poorly shielded for use in clandestine activity.
In addition, much pirate radio activity is broadcast in the FM range, because of the band's greater clarity and listenership, the smaller size and lower cost of equipment.
| Technology | Broadcasting | null |
1608705 | https://en.wikipedia.org/wiki/African%20spurred%20tortoise | African spurred tortoise | The African spurred tortoise (Centrochelys sulcata), also called the sulcata tortoise, is an endangered species of tortoise inhabiting the southern edge of the Sahara Desert, the Sahel, in Africa. It is the largest mainland species of tortoise in Africa, and the third-largest in the world, after the Galapagos tortoise and Aldabra giant tortoise. It is the only living species in its genus, Centrochelys.
Taxonomy and etymology
In 1779 the English illustrator John Frederick Miller included a hand-coloured plate of the African spurred tortoise in his Icones animalium et plantarum and coined the binomial name Testudo sulcata. Its specific name sulcata is from the Latin word sulcus meaning "furrow" and refers to the furrows on the tortoise's scales. The species was subsequently given other binomial names. There are no recognized subspecies despite there being two separate populations, one in Western Africa and the other in Eastern Africa. There are also three different, yet similar, haplotypes. One haplotype is found in and closely around Sudan, another is found in the western portion of their range, and the last haplotype is found in Senegal, Mali, and Sudan.
Range and habitat
The African spurred tortoise is native to the Sahara Desert and the Sahel, a transitional ecoregion of semiarid grasslands, savannas, and thorn shrublands found in the countries of Burkina Faso, the Central African Republic, Chad, Eritrea, Ethiopia, Mali, Mauritania, Nigeria, Senegal, Saudi Arabia, Sudan, Yemen and possibly in Somalia, Algeria, Benin, and Cameroon. It is possibly extirpated from Djibouti and Togo. They are found on hills, stable dunes, and flat areas with shrubs and high grass. They also like to settle in areas with interrupted streams or rivers. In these arid regions, the tortoise excavates burrows in the ground to get to areas with higher moisture levels, and spends the hottest part of the day in these burrows. This is known as aestivation. In the wild, they may burrow very deep, up to 15 m deep and 30 m long. Plants such as grasses and succulents grow around their burrows if kept moist, and in nature they continue to grow for the tortoise to eat if the soil is replenished with its feces. Sulcata tortoises found in the Sudanese part of their range may reach significantly greater size at maturity than those found in other regions.
Size and lifespan
C. sulcata is the largest species of tortoise in Africa and is also third-largest species of tortoise in the world. The species is the largest of the mainland tortoises. Males have an average mass of about 81 kg, but some males have been recorded at over 100 kg, with one weighing more than 120 kg. They have a straight carapace length of around 86.0 cm in males. Females have a straight carapace length of about 57.8 cm. Males of a curved carapace length of about 101.0 cm and females have approximately 67.0 cm of curved carapace length. Despite being the largest tortoise in Africa, hatchlings measure merely about 44 millimeters and weigh around 40 grams. They grow very quickly, reaching 6–10 in (15–25 cm) within the first few years of their lives. The tortoises grow faster when there is more rainfall and slower when there is less. They reach sexual maturity after 10 to 15 years. In captivity their life span is around 54 years. In the wild their lifespan is unknown but is believed to exceed 75 years. The tortoise has no known predators when they are hatchlings or adults. In fact it is believed that they are nearly immune to predators when their weight exceeds 30 kg. On the other hand tortoise eggs are sought after by many predators such as numerous species of lizards and potentially mongooses. In the wild the leading cause of death is being unable to right themselves around after they have been flipped onto their backs.
Diet
Sulcata tortoises are mostly herbivores. Primarily, their diets consist of many types of grasses, plants (especially succulent plants), and hay. Their overall diet should be high in fiber and very low in protein. Too much protein will lead to the tortoise growing too fast, which can result in metabolic bone disease, a condition that is characterized by distortion of the skeleton and weakened bone structure and can lead to lameness, lower quality of life, and/or shortened lifespan. Flowers and other plants including cactus pads can be consumed. In the wild, they have been observed to also eat plants and algae off the surface of the water. African spurred tortoises are also capable of eating various vegetables such as endive, dandelion greens, and dark leafy greens. Despite being herbivores, they will occasionally eat the carcasses of dead animals. They mostly eat dead goats and zebras that have been pushed downstream during the wet season on the rivers and streams next to which tortoises settle. If a human settlement is nearby they will also feed on refuse.
Breeding
Male selection
Copulation takes place right after the rainy season, during the months from September through November with breeding actions occur in the morning. Male C. sulcata are extremely territorial. Males combat each other for breeding rights with the females and are vocal during copulation. Larger males tend to always win sexual combat.
Female nesting
Sixty days after mating, the female begins to roam looking for suitable nesting sites. For five to 15 days, four or five nests may be excavated before she selects the perfect location in which the eggs will be laid. Females tend to lay around two to three clutches of eggs with each clutch containing 14 to 40 eggs.
Loose soil is kicked out of the depression, and the female may frequently urinate into the depression. Once it reaches about 2 feet (60 cm) in diameter and 3–6 in (7–14 cm) deep, a further depression, measuring some 8 in (20 cm) across and in depth, will be dug out towards the back of the original depression. The work of digging the nest may take up to five hours; the speed with which it is dug seems to be dependent upon the relative hardness of the ground. It usually takes place when the ambient air temperature is at least 78 °F (27 °C). Once the nest is dug, the female begins to lay an egg every three minutes. Clutches may contain 15–30 or more eggs. After the eggs are laid, the female fills in the nest, taking an hour or more to fully cover them all. Incubation should be 86 to 88 °F, and will take from 90 to 120 days.
Conservation status and efforts
Status
C. sulcata is currently ranked as an endangered species. Studies suggest that African spurred tortoises exist in approximately 16.7% of the area where they had previously been found. These studies also show an average of 1-5 tortoises per site canvassed which indicates a rapid decline of the species. The species faces threats from livestock as they have to compete for resources. The main source of resource competition African spurred tortoises face is from cattle which also graze on grass. The effects of competition for grazing land is compounded by wildfires which can destroy large portions of grass land which kills and reduces the resources available to C sulcata. They also face threats from the pet trade as they are over harvested from their natural environment. Approximately 9000 tortoises are taken from the wild for the pet trade. Other threats that the species face are habitat loss due to climate change and predators which hunt the tortoises or their eggs.
Efforts
The main method of conservation has been reintroduction programs. These sorts of reintroduction programs have been implemented in Ferlo, and Senegal. These programs have seen tortoise survival rates of about 80%. This means that the tortoises are able to easily to adapt back into their native savanna environments from domestic environment. There are also captive colonies in several countries. Most of these reintroduction programs and captive colonies can be found in protected national parks and wildlife sanctuaries. There are hopes to expand reintroduction programs by involving tortoise owners since there are more African spurred tortoises living in captivity than in the wild. The goal would be to establish a breeding program with the owners where all hatchlings would be reintroduced.
Life in captivity
Behavior
African spurred tortoises are passive and docile pets. They are almost never aggressive and barely ever show territorial behavior. This docile behavior is complemented by their slow speed and silence. Despite their docile attitude, the tortoise should not be handled often as handling will cause stress. Stress can lead to health problems and premature death. Being off the ground or constantly carried can cause it to become stressed. Stressed tortoises may urinate or defecate, try to get away, hide in its shell and hiss. They are also very curious, and can end up stuck on their backs, needing help getting flipped back over. African Spurred Tortoises have a lifespan of around 70 years, which means that they are a long term commitment.
Requirements
The ideal enclosure for the African spurred tortoise is an outdoor pen where they will be able to construct a burrow. A fence of about 2 feet in height is recommended with some parts of the fence being extended underground. They prefer high temperatures and thrive in temperatures as high as 100 degrees Fahrenheit when they have a burrow to go into to cool off. When in captivity they should also have access to heating systems to keep the temperature of an enclosure above 60 degrees Fahrenheit for when the temperature drops during the night. If the tortoise is being kept inside they need access to artificial means of sunlight. The enclosure should also be kept somewhat humid. Humidity should be kept around 40-50 percent because higher humidity may cause respiratory issues. These tend to be fungal infections, but shell rot is also common.
They require high-fiber diets (grasses and hays) as many "wet" vegetables can cause health problems in large quantities. Red leaf lettuce, prickly pear cactus pads, hibiscus leaves, hay from various grasses and dandelions are some of the better foods to make up the bulk of their diet. They will attempt to eat most types of plants eventually and some common garden plants can be very toxic to them, such as azaleas. They will eat such things as caterpillars and snails if given the opportunity, but this also should be a very small portion of their diet. Calcium should also be another small portion of their diet to help with shell growth. The tortoises should also avoid proteins and consume fruits very sparsely. As the tortoises get older and their jaws stronger, it is recommended to allow them to eat hays such as Orchard and Timothy Hay. Certain vegetables can lead to serious medical issues. Parsley, broccoli, kale and spinach should be excluded from their diets entirely, as they are too high in calcium oxalate.
| Biology and health sciences | Turtles | Animals |
1609224 | https://en.wikipedia.org/wiki/Feasibility%20study | Feasibility study | A feasibility study is an assessment of the practicality of a project or system. A feasibility study aims to objectively and rationally uncover the strengths and weaknesses of an existing business or proposed venture, opportunities and threats present in the natural environment, the resources required to carry through, and ultimately the prospects for success. In its simplest terms, the two criteria to judge feasibility are cost required and value to be attained.
A well-designed feasibility study should provide a historical background of the business or project, a description of the product or service, accounting statements, details of the operations and management, marketing research and policies, financial data, legal requirements and tax obligations. Generally, feasibility studies precede technical development and project implementation. A feasibility study evaluates the project's potential for success; therefore, perceived objectivity is an important factor in the credibility of the study for potential investors and lending institutions. It must therefore be conducted with an objective, unbiased approach to provide information upon which decisions can be based.
Formal definition
A project feasibility study is a comprehensive report that examines in detail the five frames of analysis of a given project. It also takes into consideration its four Ps, its risks and POVs, and its constraints (calendar, costs, and norms of quality). The goal is to determine whether the project should go ahead, be redesigned, or else abandoned altogether.
The five frames of analysis are:
The frame of definition;
the frame of contextual risks;
the frame of potentiality;
the parametric frame;
the frame of dominant and contingency strategies.
The four Ps are traditionally defined as Plan, Processes, People, and Power. The risks are considered to be external to the project (e.g., weather conditions) and are divided in eight categories: (Plan) financial and organizational (e.g., government structure for a private project); (Processes) environmental and technological; (People) marketing and sociocultural; and (Power) legal and political. POVs are Points of Vulnerability: they differ from risks in the sense that they are internal to the project and can be controlled or else eliminated.
The constraints are the standard constraints of calendar, costs and norms of quality that can each be objectively determined and measured along the entire project lifecycle. Depending on projects, portions of the study may suffice to produce a feasibility study; smaller projects, for example, may not require an exhaustive environmental assessment.
Common factors
TELOS is an acronym in project management used to define five areas of feasibility that determine whether a project should run or not.
T - Technical — Is the project technically possible?
E - Economic — Can the project be afforded? Will it increase profit?
L - Legal — Is the project legal?
O - Operational — How will the current operations support the change?
S - Scheduling — Can the project be done in time?
Technical feasibility
This assessment is based on an outline design of system requirements, to determine whether the company has the technical expertise to handle completion of the project. When writing a feasibility report, the following should be taken to consideration:
A brief description of the business to assess more possible factors which could affect the study
The part of the business being examined
The human and economic factor
The possible solutions to the problem
At this level, the concern is whether the proposal is both technically and legally feasible (assuming moderate cost).
The technical feasibility assessment is focused on gaining an understanding of the present technical resources of the organization and their applicability to the expected needs of the proposed system. It is an evaluation of the hardware and software and how it meets the need of the proposed system
Method of production
The selection among a number of methods to produce the same commodity should be undertaken first. Factors that make one method being preferred to other method in agricultural projects are the following:
Availability of inputs or raw materials and their quality and prices.
Availability of markets for outputs of each method and the expected prices for these outputs.
Various efficiency factors such as the expected increase in one additional unit of fertilizer or productivity of a specified crop per one thing
Production technique
After we determine the appropriate method of production of a commodity, it is necessary to look for the optimal technique to produce this commodity.
Project requirements
Once the method of production and its technique are determined, technical people have to determine the projects' requirements during the investment and operating periods. These include:
Determination of tools and equipment needed for the project such as drinkers and feeders or pumps or pipes ...etc.
Determination of projects' requirements of constructions such as buildings, storage, and roads ...etc. in addition to internal designs for these requirements.
Determination of projects' requirements of skilled and unskilled labor and managerial and financial labor.
Determination of construction period concerning the costs of designs and consultations and the costs of constructions and other tools.
Determination of minimum storage of inputs, cash money to cope with operating and contingency costs.
Project location
The most important factors that determine the selection of project location are the following:
Availability of land (proper acreage and reasonable costs).
The impact of the project on the environment and the approval of the concerned institutions for license.
The costs of transporting inputs and outputs to the project's location (i.e., the distance from the markets).
Availability of various services related to the project such as availability of extension services or veterinary or water or electricity or good roads ...etc.
Legal feasibility
It determines whether the proposed system conflicts with legal requirements, e.g., a data processing system must comply with the local data protection regulations and if the proposed venture is acceptable in accordance to the laws of the land.
Operational feasibility study
Operational feasibility is the measure of how well a proposed system solves problems and takes advantage of the opportunities identified during scope definition and how it satisfies the requirements identified in the requirements analysis phase of system development.
The operational feasibility assessment focuses on the degree to which the proposed development project fits in with the existing business environment and objectives about the development schedule, delivery date, corporate culture and existing business processes.
To ensure success, desired operational outcomes must be imparted during design and development. These include such design-dependent parameters as reliability, maintainability, supportability, usability, producibility, disposability, sustainability, affordability, etc. These parameters are required to be considered at the early stages of the design if desired operational behaviours are to be realised. A system design and development requires appropriate and timely application of engineering and management efforts to meet the previously mentioned parameters. A system may serve its intended purpose most effectively when its technical and operating characteristics are engineered into the design. Therefore, operational feasibility is a critical aspect of systems engineering that must be integral to the early design phases.
Time feasibility
A time feasibility study will take into account the period in which the project is going to take up to its completion. A project will fail if it takes too long to be completed before it is useful. Typically this means estimating how long the system will take to develop, and if it can be completed in a given time period using some methods like payback period. Time feasibility is a measure of how reasonable the project timetable is. Given our technical expertise, are the project deadlines reasonable? Some projects are initiated with specific deadlines. It is necessary to determine whether the deadlines are mandatory or desirable.
Other feasibility factors
Resource feasibility
Describe how much time is available to build the new system, when it can be built, whether it interferes with normal business operations, type and amount of resources required, dependencies, and developmental procedures with company revenue prospectus.
Financial feasibility
In case of a new project, financial viability can be judged on the following parameters:
Total estimated cost of the project
Financing of the project in terms of its capital structure, debt to equity ratio and promoter's share of total cost
Existing investment by the promoter in any other business
Projected cash flow and profitability
The financial viability of a project should provide the following information:
Full details of the assets to be financed and how liquid those assets are.
Rate of conversion to cash-liquidity (i.e., how easily the various assets can be converted to cash).
Project's funding potential and repayment terms.
Sensitivity in the repayments capability to the following factors:
Mild slowing of sales.
Acute reduction/slowing of sales.
Small increase in cost.
Large increase in cost.
Adverse economic conditions.
In 1983 the first generation of the Computer Model for Feasibility Analysis and Reporting (COMFAR), a computation tool for financial analysis of investments, was released. Since then, this United Nations Industrial Development Organization (UNIDO) software has been developed to also support the economic appraisal of projects.
The COMFAR III Expert is intended as an aid in the analysis of investment projects. The main module of the program accepts financial and economic data, produces financial and economic statements and graphical displays and calculates measures of performance. Supplementary modules assist in the analytical process. Cost-benefit and value-added methods of economic analysis developed by UNIDO are included in the program and the methods of major international development institutions are accommodated. The program is applicable for the analysis of investment in new projects and expansion or rehabilitation of existing enterprises as, e.g., in the case of reprivatisation projects. For joint ventures, the financial perspective of each partner or class of shareholder can be developed. Analysis can be performed under a variety of assumptions concerning inflation, currency revaluation and price escalations.
Market research
Market research studies is one of the most important sections of the feasibility study as it examines the marketability of the product or service and convinces readers that there is a potential market for the product or service. If a significant market for the product or services cannot be established, then there is no project. Typically, market studies will assess the potential sales of the product, absorption and market capture rates and the project's timing. The feasibility study outputs the feasibility study report, a report detailing the evaluation criteria, the study findings, and the recommendations.
| Technology | Basics | null |
1609557 | https://en.wikipedia.org/wiki/Entelodontidae | Entelodontidae | Entelodontidae is an extinct family of pig-like artiodactyls (even-toed ungulates) which inhabited the Northern Hemisphere (Asia, Europe, and North America) from the late Eocene to the early Miocene epochs, about 38-19 million years ago. Their large heads, low snouts, narrow gait, and proposed omnivorous diet inspires comparisons to suids (true pigs) and tayassuids (peccaries), and historically they have been considered closely related to these families purely on a morphological basis. However, studies which combine morphological and molecular (genetic) data on artiodactyls instead suggest that entelodonts are cetancodontamorphs, more closely related to hippos and cetaceans through their resemblance to Pakicetus, than to basal pigs like Kubanochoerus and other ungulates.
Description
Entelodonts could get quite large, and in many cases are the largest mammals in their respective ecosystems. The largest entelodont known from a complete skeleton was Daeodon, a North American entelodont which could reach an estimated weight of 750 kg (1650 pounds), and a height up to tall at the shoulder. Paraentelodon intermedium, a Eurasian species known mostly by the teeth and jaws, was similar in size to Daeodon.
Skull
Entelodonts had huge heads, ornamented with distinctive bony expansions. The zygomatic arches (cheekbones) develop huge jugal flanges which project downwards and outwards. Moreover, the underside of the lower jaw typically has one or two pairs of knob-like mandibular tubercles. These are not always diagnostic to specific taxa: often the size and presence of tubercles is variable within a single species.
The snout was narrow and elongated, especially in later species. The cranium was robust, with strong zygomatic and postorbital arches forming the rim of voluminous temporal fossae, separated by a sharp sagittal crest. However, the rear of the skull was also much shorter than the snout, and the braincase was relatively small. Most of the braincase contributed to large paranasal sinuses and olfactory bulbs at the front, while the cerebrum was underdeveloped. Large olfactory bulbs are likely indicative of a good sense of smell. Moreover, the orbits (eye sockets) are oriented further forwards than in most artiodactyls, suggesting that entelodonts had binocular vision. Compared to other artiodactyls, the jaw was slender at the rear, with a short, triangular coronoid process which is shifted forwards. The mandibular condyle (jaw joint) is set back and below the level of the tooth row. The mandibular symphysis (chin) was fused, and the pterygoid bones along the middle of the roof of the mouth were connected by a strong interdigitating suture.
Teeth
Similar to pigs, entelodonts retain a large number of teeth, a plesiomorphic trait approximating the ancestral condition for artiodactyls. They have a typical mammalian dental formula of 3.1.4.3 / 3.1.4.3, meaning that each tooth row has three pairs of robust incisors, a pair of large canines, four pairs of pointed premolars, and three pairs of relatively simple and flat molars. This unreduced, or "complete" dentition is the origin of the family's name, which is Greek for "complete teeth".
The incisors are closely packed but do not develop a distinct straight chopping surface. They range from chisel-shaped in some entelodonts (Archaeotherium) to massive and rounded in others (Daeodon). The canines have thick enamel and are circular in cross section, unlike most artiodactyls. In older individuals, the tip of the upper canine often heavily worn or even chipped off. Premolars are triangular when seen from the side, with a large and conical main cusp. They are elongated from front-to-back and widely-spaced, taking up a large portion of the tooth row. The molar teeth are bunodont, with very low and rounded cusps rather than shearing surfaces. Bunodont teeth are common in other omnivorous mammals, including pigs, bears, and humans. The upper molars have up to six cusps and a low crest (a precingulum) on the front edge of the crown. In all but the earliest entelodonts, the lower molars have only four main cusps. The front two cusps (the metaconid and protoconid) may be connected by a horizontal crest and are slightly larger than the rear two cusps.
Postcranial skeleton
The skeleton is fairly unspecialized in entelodonts. They retain typical artiodactyl skeletal traits such as a double-pulley ankle joint and paraxonic ("even toed") feet with weight split evenly between the two middle toes. They had four toes in total, with the middle two forming small, pointed cloven hooves, while the remaining two were vestigial and likely not externally visible. In larger species, a bison-like spinal hump supported the weight of the heavy head. The limbs were long, and the radius and ulna were fused. Though not fused, the metatarsals (raised foot bones) were long and closely packed. The limb and hoof proportions are consistent with other hoofed animals that run well on open ground but are not built for high speed.
Paleobiology
Jaw movement and musculature
The wide and tall temporal fossa allowed for a very large temporalis muscle, which extends from the side of the cranium to the coronoid process of the mandible. The temporalis was not only large and strong, but also had a long moment arm (and thus higher torque) due to the coronoid process shifting forwards. The reinforced pterygoid, zygomatic, and postorbital areas would have supported the forces generated by the temporalis. The size and orientation of the temporalis is similar to carnivorans, where it corresponds to a strong and stable scissor-like (orthal) bite.
Though the low jaw joint provided more room for the temporalis muscle, it also posed a problem for the masseter muscle. The masseter, which extends from the zygomatic arch to the lower rear corner of the mandible, is a major component of the chewing apparatus in herbivorous artiodactyls. While other artiodactyls added torque to the muscle by raising the jaw joint, entelodonts instead expanded the rear of the jaw downwards, as a deep, curved flange. Moreover, the characteristic jugal flanges of entelodonts were covered with muscle scars on the inside, likely attachment points to strengthen the masseter. Only a few modern mammals have overdeveloped projections on the zygomatic arch, including xenarthrans, kangaroos, and certain rodents. Like entelodonts, these mammals use their equivalent projections as a means of providing extra space for the attachment of the masseter muscle, and develop robust cranial bars to resist the resulting forces on the skull. The pterygoideus muscle, which follows a similar path and function to the masseter, also benefited from the deep flange at the back of the jaw.
The function of the mandibular tubercles is not certain, but they may also be related to jaw musculature. They are only clearly correlated with the size of the individual, though a few taxa (Brachyhyops and Cypretherium) can be diagnosed by the absence of a specific pair of mandibular tubercles. Generally, the posterior (rear) mandibular tubercles develop later in life than the anterior (front) pair, and none of the tubercles stop growing as the animal develops. The use of the anterior tubercles is unclear; one speculative idea suggests that they served as an attachment point for strong lip muscles in particularly herbivorous entelodonts. The posterior tubercles may provide a link to the digastricus muscle which helps to open the jaws. Hippos, which have a particularly complex and well-developed digastricus, occasionally develop a tubercle to support the digastricus in an equivalent area on the jaw.
The jaw joint of entelodonts was likely more strongly connected than the loose jaws of most other artiodactyls. The mandibular condyle was convex and inserted into a strongly concave facet (glenoid) on the zygomatic arch, which would have restricted front-to-back (propalinal) jaw movement. Nevertheless, the structure of the mandibular condyle itself allowed for a wide range of movement, and the laterally bowed zygomatic arch provided some room for side-to-side (transverse) movement driven by the masseter and pterygoideus. The low, unconstrained jaw joint and short coronoid process may correspond to long muscle fibers. This points to a hinge-like jaw suspension with a very wide gape, similar to some modern carnivorans such as felids (cats). Based on the shape of the mandibular condyle, the maximum gape possible based on the underlying bones (though not necessarily the widest gape possible in life) was about 109 degrees in Archaeotherium.
Wear facets on entelodont teeth support three-part food processing. First, the incisors and canines bite in a strong orthal motion, grabbing and puncturing food. Then, the food is transferred back to the premolars, which breaks apart tough parts of the food with similar movements. Finally, the food is crushed and ground up by the molars, using a combination of orthal and transverse grinding. This same basic process is seen in modern pigs and peccaries, which have similar dentition. Individuals may have preferred one side of the jaw for chewing, as premolars and molars often show an asymmetrical distribution of wear between the left and right sides of the mouth.
Diet
By comparison to pigs and peccaries, entelodonts were almost certainly omnivorous to an extent. Their teeth and jaw structure would have assisted processing of large and tough food items. Unlike the diverse and fully herbivorous pecoran artiodactyls, entelodonts lack specializations for chopping and shredding grass and other particularly fibrous plants. Instead, entelodonts were probably browsers, with roots, nuts, fruits, and branches as their preferred sources of vegetation. A 2022 study found that Entelodon magnus had an omnivorous diet similar to wild boar (Sus scrofa). This conclusion was justified by its pattern of tooth microwear, run through a linear discriminant analysis calibrated by modern herbivorous and omnivorous mammals. Based on pigs, entelodonts probably had a simple stomach and relied on the caecum to ferment and digest plant matter. They would have been opportunistic omnivores, capable of digesting a variety of plant and animal matter and moderating their food preferences based on seasonal ability.
The same adaptations useful for processing tough plant material would be equally useful for carrion and bones, which could have been major components of the diet for some entelodonts. Unlike pigs, the youngest juvenile entelodonts had a full set of 32 deciduous teeth. The teeth were sharp, slender, and semi-serrated, less suitable for crushing tough food compared to adult entelodonts.
In many entelodonts, the canine teeth acquire rounded wear surfaces at their tips, indicating regular use on hard material such as bones. Similar patterns of canine wear are observed in modern cats, which rely on strong bites administered through their canine teeth when killing prey. In some species the bases of the canines are scoured by smooth grooves, a trait consistent with abrasions from sediment-covered plant material such as roots. These grooves instead could have been produced by stripping long, fibrous vegetation, such as water-rich grape vines. Daeodon is known to have had a distinctive type of "piecrust" tooth wear at the tips of the premolars, with a flat dentine surface surrounded by chipped enamel. This has also been observed in living hyenas. Few contemporary mammals approached entelodonts in the extent of adaptations consistent with scavenging. Fossils with large scrapes and puncture marks are found throughout entelodont-bearing sites in the American Great Plains, including a skull of Merycoidodon with an embedded incisor of the entelodont Archaeotherium.
Entelodonts may have engaged in active predation, though the extent of this behavior is debated. Several species of modern pigs occasionally engage in predation, and even traditional herbivores like camels show dental wear consistent with scavenging. If they did engage in predation, entelodonts would not have been alone: many other contemporary mammals filled apex predator niches, including cat-like saber-toothed nimravids, amphicyonids ("bear-dogs"), and hyaenodontid creodonts. One of the most apparent examples of circumstantial evidence for predation is a fossil found in the White River Formation of Wyoming, representing a cache of partial skeletons and other remains of the early camelid Poebrotherium. The carcasses were covered with large punctures on the skull, neck, and the transition from the thoracic to lumbar vertebrae, which have been attributed to predation and scavenging by Archaeotherium. Entelodon's tooth microwear showed no overlap with the modern brown bear (Ursus arctos), and it probably did not actively hunt large mammals as part of its normal diet.
Intraspecific behavior
The jaw structure and estimated musculature hold numerous lines of evidence indicating that entelodonts could open their mouths unusually wide. This trait may have been useful in hunting or feeding on carrion, but similar adaptations have also been linked to competitive behaviors in herbivores. Hippos, a related group with similar adaptations, are aggressive herbivores which can open their jaws up to 150 degrees and display enlarged canines in order to intimidate rivals. Male hippos engage in head-to-head "yawning" and jaw-wrestling contests, while females attack by approaching from the side and slamming their head into the opponent's body. The wide gape and low skulls of entelodonts would have assisted biting competitions, which are supported by fossil evidence. Large bite marks, including healed punctures, are common on skulls of various American entelodonts. These wounds are concentrated above the sinuses, and are only found on adult specimens. One could easily draw comparisons between these bite marks and the wide range of intraspecific competition over mates or territories in modern artiodactyls. Snout biting in particular is a common competitive behavior among male camels, another group of "primitive" artiodactyls. Ribcage injuries have been attributed to intraspecies aggression in Archaeotherium. One possible function for the anterior tubercles is as a support for toughened skin, which would have acted as a buffer or display feature during competitive behavior.
Classification
Early history
The earliest entelodont fossils to be named were described within a short time frame in the 1840s. The first entelodont species known from good fossils was Entelodon magnus, a European species which was named by French paleontologist Auguste Aymard. There is some debate over when Aymard's description was first published; though most authors assumed it was written in 1846, a citation within the article suggests that it was not published until 1848. Auguste Pomel, one of Aymard's contemporaries, described another fossil as Elotherium around the same time. Pomel's volume was likely published in 1846 or 1847, albeit with surviving reprints dating to 1848. Entelodon and Elotherium are almost certainly synonymous, though fossils belonging to the latter name are fragmentary and have been lost, while those of the former were likely described later. Nearly all historical and modern authors prefer to use Entelodon for the purpose of clarity, even though it would not take priority under strict rules of nomenclature. The confusion of priority between Entelodon and Elotherium is reflected in the name of their corresponding family. Edward Richard Alston coined the name Elotheriidae in 1878, while Richard Lydekker used the name Entelodontidae in 1883. As with Entelodon, nearly all paleontologists prefer Entelodontidae when referring to the family.
Following the confusion between Entelodon and Elotherium, entelodont fossils continued to be discovered in Europe. Large entelodonts were also described from North America starting in 1850, though most new genera were eventually lumped into Archaeotherium and Daeodon. By the beginning of the 20th century, entelodont skeletal anatomy was well-understood from the quantity of fossils discovered by that point. In 1909, a massive complete skeleton of "Dinohyus" hollandi (= Daeodon), CM 1594, was described and put on display at the Carnegie Museum of Natural History. As the 20th century continued, Asian entelodonts were discovered (Eoentelodon, Paraentelodon), as well as some of the earliest known members of the family (Eoentelodon, Brachyhyops).
Traditional classification
The first described entelodonts were described in conjunction with Richard Owen's recognition of the artiodactyls as a natural group. The earliest sources considered entelodonts to be true pigs, but as further fossils were discovered, it became clear that they had a long evolutionary history separate from pigs. Regardless, entelodonts were universally accepted as examples of "primitive" artiodactyls, with unspecialized bunodont teeth in contrast with the strong adaptations for herbivory present in the more "advanced" ruminants. Various names were erected to encompass living and extinct bunodont-toothed and non-ruminant artiodactyls, such as "Omnivoria" (Owens, 1858), "Bunodontia" (Lydekker, 1883) and "Nonruminantia" (Gregory, 1910).
Some authors considered entelodonts to be too "primitive" for comparison to modern bunodont artiodactyls. In these studies, entelodonts were placed in "Palaeodonta", a group shared with various other extinct families. Choeropotamids, cebochoerids, and helohyids were frequently associated with entelodonts, sometimes even as potential ancestors. Later, the superfamily Entelodontoidea was named to encompass Entelodontidae and their supposed closest extinct relatives. In modern studies, Entelodontidae is generally considered the only family within Entelodontoidea.
Many studies argued that entelodonts had close relations to living pigs, peccaries, and hippos. Various groups have been developed and named in reference to a pig-like anatomy, with names such as Suina (Gray, 1868) and Suiformes (Jaeckel, 1911) being emplaced in varying contexts. A restricted definition of Suina is still in use, as a major artiodactyl suborder encompassing Tayassuidae (peccaries) and Suidae (pigs). Early cladistic phylogenetic analyses of artiodactyls placed Entelodontidae as the sister taxon to a Tayassuidae + Suidae clade. This seemed to justify the frequent morphological comparisons between entelodonts and pigs.
Cetancodontomorpha
While entelodonts have long been classified as members of the Suina, Spaulding et al. have found them to be closer to whales and hippos than to pigs. Cladistic analysis of the position of whales in relation to artiodactyls and mesonychians changes radically depending on whether the giant enigmatic mammal Andrewsarchus is included, and it has been suggested that Andrewsarchus is in fact an entelodont or close relative.
Many former genera of entelodonts have been synonymized. For example, some authors have synonymized Dinohyus with Daeodon shoshonensis, a species described from fragmentary material by Cope.
List of genera
†Archaeotherium
†Brachyhyops
†Cypretherium
†Daeodon
†Entelodon
†Entelodontellus
†Eoentelodon
†Paraentelodon
†Proentelodon? (may not be an entelodont)
In popular culture
In popular media, entelodonts are sometimes nicknamed hell pigs or terminator pigs.
Entelodonts appear in the third episode of the popular BBC documentary Walking with Beasts, where, in the program, the narrator always refers to the creatures as "entelodonts" rather than a more specific genus, such as Entelodon. The same creatures appear in another BBC production, the 2001 remake of The Lost World.
Entelodonts were also the main focus of episode 4 of National Geographic Channel's show Prehistoric Predators in an episode titled "Killer Pig". The episode featured a number of claims unproven or disproven by science, such as Archaeotherium (identified as "entelodont") being the top predator of the American Badlands, and evolving directly into the even larger Daeodon (called "Dinohyus" in the episode).
| Biology and health sciences | Other artiodactyla | Animals |
1610231 | https://en.wikipedia.org/wiki/Energy%20density | Energy density | In physics, energy density is the quotient between the amount of energy stored in a given system or contained in a given region of space and the volume of the system or region considered. Often only the useful or extractable energy is measured. It is sometimes confused with stored energy per unit mass, which is called specific energy or .
There are different types of energy stored, corresponding to a particular type of reaction. In order of the typical magnitude of the energy stored, examples of reactions are: nuclear, chemical (including electrochemical), electrical, pressure, material deformation or in electromagnetic fields. Nuclear reactions take place in stars and nuclear power plants, both of which derive energy from the binding energy of nuclei. Chemical reactions are used by organisms to derive energy from food and by automobiles from the combustion of gasoline. Liquid hydrocarbons (fuels such as gasoline, diesel and kerosene) are today the densest way known to economically store and transport chemical energy at a large scale (1 kg of diesel fuel burns with the oxygen contained in ≈ 15 kg of air). Burning local biomass fuels supplies household energy needs (cooking fires, oil lamps, etc.) worldwide. Electrochemical reactions are used by devices such as laptop computers and mobile phones to release energy from batteries.
Energy per unit volume has the same physical units as pressure, and in many situations is synonymous. For example, the energy density of a magnetic field may be expressed as and behaves like a physical pressure. The energy required to compress a gas to a certain volume may be determined by multiplying the difference between the gas pressure and the external pressure by the change in volume. A pressure gradient describes the potential to perform work on the surroundings by converting internal energy to work until equilibrium is reached.
In cosmological and other contexts in general relativity, the energy densities considered relate to the elements of the stress–energy tensor and therefore do include the rest mass energy as well as energy densities associated with pressure.
Chemical energy
When discussing the chemical energy contained, there are different types which can be quantified depending on the intended purpose. One is the theoretical total amount of thermodynamic work that can be derived from a system, at a given temperature and pressure imposed by the surroundings, called exergy. Another is the theoretical amount of electrical energy that can be derived from reactants that are at room temperature and atmospheric pressure. This is given by the change in standard Gibbs free energy. But as a source of heat or for use in a heat engine, the relevant quantity is the change in standard enthalpy or the heat of combustion.
There are two kinds of heat of combustion:
The higher value (HHV), or gross heat of combustion, includes all the heat released as the products cool to room temperature and whatever water vapor is present condenses.
The lower value (LHV), or net heat of combustion, does not include the heat which could be released by condensing water vapor, and may not include the heat released on cooling all the way down to room temperature.
A convenient table of HHV and LHV of some fuels can be found in the references.
In energy storage and fuels
For energy storage, the energy density relates the stored energy to the volume of the storage equipment, e.g. the fuel tank. The higher the energy density of the fuel, the more energy may be stored or transported for the same amount of volume. The energy of a fuel per unit mass is called its specific energy.
The adjacent figure shows the gravimetric and volumetric energy density of some fuels and storage technologies (modified from the Gasoline article). Some values may not be precise because of isomers or other irregularities. The heating values of the fuel describe their specific energies more comprehensively.
The density values for chemical fuels do not include the weight of the oxygen required for combustion. The atomic weights of carbon and oxygen are similar, while hydrogen is much lighter. Figures are presented in this way for those fuels where in practice air would only be drawn in locally to the burner. This explains the apparently lower energy density of materials that contain their own oxidizer (such as gunpowder and TNT), where the mass of the oxidizer in effect adds weight, and absorbs some of the energy of combustion to dissociate and liberate oxygen to continue the reaction. This also explains some apparent anomalies, such as the energy density of a sandwich appearing to be higher than that of a stick of dynamite.
Given the high energy density of gasoline, the exploration of alternative media to store the energy of powering a car, such as hydrogen or battery, is strongly limited by the energy density of the alternative medium. The same mass of lithium-ion storage, for example, would result in a car with only 2% the range of its gasoline counterpart. If sacrificing the range is undesirable, much more storage volume is necessary. Alternative options are discussed for energy storage to increase energy density and decrease charging time, such as supercapacitors.
No single energy storage method boasts the best in specific power, specific energy, and energy density. Peukert's law describes how the amount of useful energy that can be obtained (for a lead-acid cell) depends on how quickly it is pulled out.
Efficiency
In general an engine will generate less kinetic energy due to inefficiencies and thermodynamic considerations—hence the specific fuel consumption of an engine will always be greater than its rate of production of the kinetic energy of motion.
Energy density differs from energy conversion efficiency (net output per input) or embodied energy (the energy output costs to provide, as harvesting, refining, distributing, and dealing with pollution all use energy). Large scale, intensive energy use impacts and is impacted by climate, waste storage, and environmental consequences.
Nuclear energy
The greatest energy source by far is matter itself, according to the mass–energy equivalence. This energy is described by , where c is the speed of light. In terms of density, , where ρ is the volumetric mass density, V is the volume occupied by the mass. This energy can be released by the processes of nuclear fission (~ 0.1%), nuclear fusion (~ 1%), or the annihilation of some or all of the matter in the volume V by matter–antimatter collisions (100%).
The most effective ways of accessing this energy, aside from antimatter, are fusion and fission. Fusion is the process by which the sun produces energy which will be available for billions of years (in the form of sunlight and heat). However as of 2024, sustained fusion power production continues to be elusive. Power from fission in nuclear power plants (using uranium and thorium) will be available for at least many decades or even centuries because of the plentiful supply of the elements on earth, though the full potential of this source can only be realized through breeder reactors, which are, apart from the BN-600 reactor, not yet used commercially.
Fission reactors
Nuclear fuels typically have volumetric energy densities at least tens of thousands of times higher than chemical fuels. A 1 inch tall uranium fuel pellet is equivalent to about 1 ton of coal, 120 gallons of crude oil, or 17,000 cubic feet of natural gas. In light-water reactors, 1 kg of natural uranium – following a corresponding enrichment and used for power generation– is equivalent to the energy content of nearly 10,000 kg of mineral oil or 14,000 kg of coal. Comparatively, coal, gas, and petroleum are the current primary energy sources in the U.S. but have a much lower energy density.
The density of thermal energy contained in the core of a light-water reactor (pressurized water reactor (PWR) or boiling water reactor (BWR)) of typically ( electrical corresponding to ≈ thermal) is in the range of 10 to 100 MW of thermal energy per cubic meter of cooling water depending on the location considered in the system (the core itself (≈ ), the reactor pressure vessel (≈ ), or the whole primary circuit (≈ )). This represents a considerable density of energy that requires a continuous water flow at high velocity at all times in order to remove heat from the core, even after an emergency shutdown of the reactor.
The incapacity to cool the cores of three BWRs at Fukushima after the 2011 tsunami and the resulting loss of external electrical power and cold source caused the meltdown of the three cores in only a few hours, even though the three reactors were correctly shut down just after the Tōhoku earthquake. This extremely high power density distinguishes nuclear power plants (NPP's) from any thermal power plants (burning coal, fuel or gas) or any chemical plants and explains the large redundancy required to permanently control the neutron reactivity and to remove the residual heat from the core of NPP's.
Antimatter–matter annihilation
Because antimatter–matter interactions result in complete conversion of the rest mass to radiant energy, the energy density of this reaction depends on the density of the matter and antimatter used. A neutron star would approximate the most dense system capable of matter-antimatter annihilation. A black hole, although denser than a neutron star, does not have an equivalent anti-particle form, but would offer the same 100% conversion rate of mass to energy in the form of Hawking radiation. Even in the case of relatively small black holes (smaller than astronomical objects) the power output would be tremendous.
Electric and magnetic fields
Electric and magnetic fields can store energy and its density relates to the strength of the fields within a given volume. This (volumetric) energy density is given by
where is the electric field, is the magnetic field, and and are the permittivity and permeability of the surroundings respectively. The SI unit is the joule per cubic metre.
In ideal (linear and nondispersive) substances, the energy density is
where is the electric displacement field and is the magnetizing field. In the case of absence of magnetic fields, by exploiting Fröhlich's relationships it is also possible to extend these equations to anisotropic and nonlinear dielectrics, as well as to calculate the correlated Helmholtz free energy and entropy densities.
In the context of magnetohydrodynamics, the physics of conductive fluids, the magnetic energy density behaves like an additional pressure that adds to the gas pressure of a plasma.
Pulsed sources
When a pulsed laser impacts a surface, the radiant exposure, i.e. the energy deposited per unit of surface, may also be called energy density or fluence.
Table of material energy densities
The following unit conversions may be helpful when considering the data in the tables: 3.6 MJ = 1 kW⋅h ≈ 1.34 hp⋅h. Since 1 J = 10−6 MJ and 1 m3 = 103 L, divide joule/m3 by 109 to get MJ/L = GJ/m3. Divide MJ/L by 3.6 to get kW⋅h/L.
Chemical reactions (oxidation)
Unless otherwise stated, the values in the following table are lower heating values for perfect combustion, not counting oxidizer mass or volume. When used to produce electricity in a fuel cell or to do work, it is the Gibbs free energy of reaction (ΔG) that sets the theoretical upper limit. If the produced is vapor, this is generally greater than the lower heat of combustion, whereas if the produced is liquid, it is generally less than the higher heat of combustion. But in the most relevant case of hydrogen, ΔG is 113 MJ/kg if water vapor is produced, and 118 MJ/kg if liquid water is produced, both being less than the lower heat of combustion (120 MJ/kg).
Electrochemical reactions (batteries)
Common battery formats
Nuclear reactions
In material deformation
The mechanical energy storage capacity, or resilience, of a Hookean material when it is deformed to the point of failure can be computed by calculating tensile strength times the maximum elongation dividing by two. The maximum elongation of a Hookean material can be computed by dividing stiffness of that material by its ultimate tensile strength. The following table lists these values computed using the Young's modulus as measure of stiffness:
Other release mechanisms
| Physical sciences | Physics basics: General | Physics |
9567916 | https://en.wikipedia.org/wiki/Mechanical%20screening | Mechanical screening | Mechanical screening, often just called screening, is the practice of taking granulated or crushed ore material and separating it into multiple grades by particle size.
This practice occurs in a variety of industries such as mining and mineral processing, agriculture, pharmaceutical, food, plastics, and recycling.
A method of separating solid particles according to size alone is called screening.
General categories
Screening falls under two general categories: dry screening, and wet screening. From these categories, screening separates a flow of material into grades, these grades are then either further processed to an intermediary product or a finished product. Additionally, the machines can be categorized into a moving screen and static screen machines, as well as by whether the screens are horizontal or inclined.
Applications
The mining and mineral processing industry uses screening for a variety of processing applications. For example, after mining the minerals, the material is transported to a primary crusher. Before crushing large boulder are scalped on a shaker with thick shielding screening. Further down stream after crushing the material can pass through screens with openings or slots that continue to become smaller. Finally, screening is used to make a final separation to produce saleable products based on a grade or a size range.
Process
A screening machine consist of a drive that induces vibration, a screen media that causes particle separation, and a deck which holds the screen media and the drive and is the mode of transport for the vibration.
There are physical factors that makes screening practical. For example, vibration, g force, bed density, and material shape all facilitate the rate or cut. Electrostatic forces can also hinder screening efficiency in way of water attraction causing sticking or plugging, or very dry material generate a charge that causes it to attract to the screen itself.
As with any industrial process there is a group of terms that identify and define what screening is. Terms like blinding, contamination, frequency, amplitude, and others describe the basic characteristics of screening, and those characteristics in turn shape the overall method of dry or wet screening.
In addition, the way a deck is vibrated differentiates screens. Different types of motion have their advantages and disadvantages. In addition media types also have their different properties that lead to advantages and disadvantages.
Finally, there are issues and problems associated with screening. Screen tearing, contamination, blinding, and dampening all affect screening efficiency.
Physical principles
Vibration - either sinusoidal vibration or gyratory vibration.
Sinusoidal Vibration occurs at an angled plane relative to the horizontal. The vibration is in a wave pattern determined by frequency and amplitude.
Gyratory Vibration occurs at near level plane at low angles in a reciprocating side to side motion.
Gravity - This physical interaction is after material is thrown from the screen causing it to fall to a lower level. Gravity also pulls the particles through the screen media.
Density - The density of the material relates to material stratification.
Electrostatic Force - This force applies to screening when particles are extremely dry or is wet.
Screening terminology
Like any mechanical and physical entity there are scientific, industrial, and layman terminology. The following is a partial list of terms that are associated with mechanical screening.
Amplitude - This is a measurement of the screen cloth as it vertically peaks to its tallest height and troughs to its lowest point. Measured in multiples of the acceleration constant g (g-force).
Acceleration - Applied Acceleration to the screen mesh in order to overcome the van der waal forces
Blinding - When material plugs into the open slots of the screen cloth and inhibits overflowing material from falling through.
Brushing - This procedure is performed by an operator who uses a brush to brush over the screen cloth to dislodged blinded opening.
Contamination - This is unwanted material in a given grade. This occurs when there is oversize or fine size material relative to the cut or grade. Another type of contamination is foreign body contamination.
Oversize contamination occurs when there is a hole in the screen such that the hole is larger than the mesh size of the screen. Other instances where oversize occurs is material overflow falling into the grade from overhead, or there is the wrong mesh size screen in place.
Fines contamination is when large sections of the screen cloth is blinded over, and material flowing over the screen does not fall through. The fines are then retained in the grade.
Foreign body contamination is unwanted material that differs from the virgin material going over and through the screen. It can be anything ranging from tree twigs, grass, metal slag to other mineral types and composition. This contamination occurs when there is a hole in the scalping screen or a foreign material's mineralogy or chemical composition differs from the virgin material.
Deck - a deck is frame or apparatus that holds the screen cloth in place. It also contains the screening drive. It can contain multiple sections as the material travels from the feed end to the discharge end. Multiple decks are screen decks placed in a configuration where there are a series of decks attached vertically and lean at the same angle as it preceding and exceeding decks. Multiple decks are often referred to as single deck, double deck, triple deck, etc.
Frequency - Measured in hertz (Hz) or revolutions per minute (RPM). Frequency is the number of times the screen cloth sinusoidally peaks and troughs within a second. As for a gyratory screening motion it is the number of revolutions the screens or screen deck takes in a time interval, such as revolution per minute (RPM).
Gradation, grading - Also called "cut" or "cutting." Given a feed material in an initial state, the material can be defined to have a particle size distribution. Grading is removing the maximum size material and minimum size material by way of mesh selection.
Screen Media (Screen cloth) - it is the material defined by mesh size, which can be made of any type of material such steel, stainless steel, rubber compounds, polyurethane, brass, etc.
Shaker - the whole assembly of any type mechanical screening machine.
Stratification - This phenomenon occurs as vibration is passed through a bed of material. This causes coarse (larger) material to rise and finer (smaller) material to descend within the bed. The material in contact with screen cloth either falls through a slot or blinds the slot or contacts the cloth material and is thrown from the cloth to fall to the next lower level.
Mesh - The number of open slots per linear inch. Mesh is arranged in multiple configuration. Mesh can be a square pattern, long-slotted rectangular pattern, circular pattern, or diamond pattern.
Scalp, scalping - this is the very first cut of the incoming material with the sum of all its grades. Scalping is removing the largest size particles. This includes enormously large particles relative to the other particle's sizes. Scalping also cleans the incoming material from foreign body contamination such as twigs, trash, glass, or other unwanted oversize material.
Types of mechanical screening
There are a number of types of mechanical screening equipment that cause segregation. These types are based on the motion of the machine through its motor drive.
Circle-throw vibrating equipment - This type of equipment has an eccentric shaft that causes the frame of the shaker to lurch at a given angle. This lurching action literally throws the material forward and up. As the machine returns to its base state the material falls by gravity to physically lower level. This type of screening is used also in mining operations for large material with sizes that range from six inches to +20 mesh.
High frequency vibrating equipment - This type of equipment drives the screen cloth only. Unlike above the frame of the equipment is fixed and only the screen vibrates. However, this equipment is similar to the above such that it still throws material off of it and allows the particles to cascade down the screen cloth. These screens are for sizes smaller than 1/8 of an inch to +150 mesh.
Gyratory equipment - This type of equipment differs from the above two such that the machine gyrates in a circular motion at a near level plane at low angles. The drive is an eccentric gear box or eccentric weights.
Trommel screens - Does not require vibration. Instead, material is fed into a horizontal rotating drum with screen panels around the diameter of the drum.
Tumbler screening technique
An improvement on vibration, vibratory, and linear screeners, a tumbler screener uses elliptical action which aids in screening of even very fine material. As like panning for gold, the fine particles tend to stay towards the center and the larger go to the outside. It allows for segregation and unloads the screen surface so that it can effectively do its job. With the addition of multiple decks and ball cleaning decks, even difficult products can be screened at high capacity to very fine separations.
Circle-throw vibrating equipment
Circle-Throw Vibrating Equipment is a shaker or a series of shakers as to where the drive causes the whole structure to move. The structure extends to a maximum throw or length and then contracts to a base state. A pattern of springs are situated below the structure to where there is vibration and shock absorption as the structure returns to the base state.
This type of equipment is used for very large particles, sizes that range from pebble size on up to boulder size material. It is also designed for high volume output. As a scalper, this shaker will allow oversize material to pass over and fall into a crusher such a cone crusher, jaw crusher, or hammer mill. The material that passes the screen by-passes the crusher and is conveyed and combined with the crush material.
Also this equipment is used in washing processes, as material passes under spray bars, finer material and foreign material is washed through the screen. This is one example of wet screening.
High frequency vibrating equipment
High-frequency vibrating screening equipment is a shaker whose frame is fixed and the drive vibrates only the screen cloth. High frequency vibration equipment is for particles that are in this particle size range of an 1/8 in (3 mm) down to a +150 mesh. Traditional shaker screeners have a difficult time making separations at sizes like 44 microns. At the same time, other high energy sieves like the Elcan Industries' advanced screening technology allow for much finer separations down to as fine as 10um and 5um, respectively.
These shakers usually make a secondary cut for further processing or make a finished product cut.
These shakers are usually set at a steep angle relative to the horizontal level plane. Angles range from 25 to 45 degrees relative to the horizontal level plane.
Gyratory equipment
This type of equipment has an eccentric drive or weights that causes the shaker to travel in an orbital path. The material rolls over the screen and falls with the induction of gravity and directional shifts. Rubber balls and trays provide an additional mechanical means to cause the material to fall through. The balls also provide a throwing action for the material to find an open slot to fall through.
The shaker is set a shallow angle relative to the horizontal level plane. Usually, no more than 2 to 5 degrees relative to the horizontal level plane.
These types of shakers are used for very clean cuts. Generally, a final material cut will not contain any oversize or any fines contamination.
These shakers are designed for the highest attainable quality at the cost of a reduced feed rate.
Trommel screens
Trommel screens have a rotating drum on a shallow angle with screen panels around the diameter of the drum. The feed material always sits at the bottom of the drum and, as the drum rotates, always comes into contact with clean screen. The oversize travels to the end of the drum as it does not pass through the screen, while the undersize passes through the screen into a launder below.
Screen Media Attachment Systems
There are many ways to install screen media into a screen box deck (shaker deck). Also, the type of attachment system has an influence on the dimensions of the media.
Tensioned screen media
Tensioned screen cloth is typically 4 feet by the width or the length of the screening machine depending on whether the deck is side or end tensioned. Screen cloth for tensioned decks can be made with hooks and are attached with clamp rails bolted on both sides of the screen box. When the clamp rail bolts are tightened, the cloth is tensioned or even stretched in the case of some types of self-cleaning screen media. To ensure that the center of the cloth does not tap repeatedly on the deck due to the vibrating shaker and that the cloth stays tensioned, support bars are positioned at different heights on the deck to create a crown curve from hook to hook on the cloth. Tensioned screen cloth is available in various materials: stainless steel, high carbon steel and oil tempered steel wires, as well as moulded rubber or polyurethane and hybrid screens (a self-cleaning screen cloth made of rubber or polyurethane and metal wires).
Commonly, vibratory-type screening equipment employs rigid, circular sieve frames to which woven wire mesh is attached. Conventional methods of producing tensioned meshed screens has given way in recent years to bonding, whereby the mesh is no longer tensioned and trapped between a sieve frame body and clamping ring; instead, developments in modern adhesive technologies has allowed the industry to adopt high strength structural adhesives to bond tensioned mesh directly to frames.
Modular screen media
Modular screen media is typically 1 foot large by 1 or 2 feet long (4 feet long for ISEPREN WS 85 ) steel reinforced polyurethane or rubber panels. They are installed on a flat deck (no crown) that normally has a larger surface than a tensioned deck. This larger surface design compensates for the fact that rubber and polyurethane modular screen media offers less open area than wire cloth. Over the years, numerous ways have been developed to attach modular panels to the screen deck stringers (girders). Some of these attachment systems have been or are currently patented. Self-cleaning screen media is also available on this modular system.
Types of Screen Media
There are several types of screen media manufactured with different types of material that use the two common types of screen media attachment systems, tensioned and modular.
Woven Wire Cloth (Mesh)
Woven wire cloth, typically produced from stainless steel, is commonly employed as a filtration medium for sieving in a wide range of industries. Most often woven with a plain weave, or a twill weave for the lightest of meshes, apertures can be produced from a few microns upwards (e.g. 25 microns), employing wires with diameters from as little as 25 microns. A twill weave allows a mesh to be woven when the wire diameter is too thick in proportion to the aperture. Other, less commonplace, weaves, such as Dutch/Hollander, allow the production of meshes that are stronger and/or having smaller apertures.
Today wire cloth is woven to strict international standards, e.g. ISO1944:1999, which dictates acceptable tolerance regarding nominal mesh count and blemishes. The nominal mesh count, to which mesh is generally defined is a measure of the number of openings per lineal inch, determined by counting the number of openings from the centre of one wire to the centre of another wire one lineal inch away. For example, a 2 mesh woven with a wire of 1.6mm wire diameter has an aperture of 11.1mm (see picture below of a 2 mesh with an intermediate crimp). The formula for calculating the aperture of a mesh, with a known mesh count and wire diameter, is as follows:
where a = aperture, b = mesh count and c = wire diameter.
Other calculations regarding woven wire cloth/mesh can be made including weight and open area determination. Of note, wire diameters are often referred to by their standard wire gauge (swg); e.g. a 1.6mm wire is a 16 swg.
Traditionally, screen cloth was made with metal wires woven with a loom. Today, woven cloth is still widely used primarily because they are less expensive than other types of screen media. Over the years, different weaving techniques have been developed; either to increase the open area percentage or add wear-life. Slotted opening woven cloth is used where product shape is not a priority and where users need a higher open area percentage. Flat-top woven cloth is used when the consumer wants to increase wear-life. On regular woven wire, the crimps (knuckles on woven wires) wear out faster than the rest of the cloth resulting in premature breakage. On flat-top woven wire, the cloth wears out equally until half of the wire diameter is worn, resulting in a longer wear life. Unfortunately flat-top woven wire cloth is not widely used because of the lack of crimps that causes a pronounced reduction of passing fines resulting in premature wear of con crushers.
Perforated & Punch Plate
On a crushing and screening plant, punch plates or perforated plates are mostly used on scalper vibrating screens, after raw products pass on grizzly bars. Most likely installed on a tensioned deck, punch plates offer excellent wear life for high-impact and high material flow applications.
Synthetic screen media (typically rubber or polyurethane)
Synthetic screen media is used where wear life is an issue. Large producers such as mines or huge quarries use them to reduce the frequency of having to stop the plant for screen deck maintenance. Rubber is also used as a very resistant high-impact screen media material used on the top deck of a scalper screen. To compete with rubber screen media fabrication, polyurethane manufacturers developed screen media with lower Shore Hardness. To compete with self-cleaning screen media that is still primarily available in tensioned cloth, synthetic screen media manufacturers also developed membrane screen panels, slotted opening panels and diamond opening panels. Due to the 7-degree demoulding angle, polyurethane screen media users can experience granulometry changes of product during the wear life of the panel.
Self-Cleaning Screen Media
Self-cleaning screen media was initially engineered to resolve screen cloth blinding, clogging and pegging problems. The idea was to place crimped wires side by side on a flat surface, creating openings and then, in some way, holding them together over the support bars (crown bars or bucker bars). This would allow the wires to be free to vibrate between the support bars, preventing blinding, clogging and pegging of the cloth. Initially, crimped longitudinal wires on self-cleaning cloth were held together over support bars with woven wire. In the 50s, some manufacturers started to cover the woven cross wires with caulking or rubber to prevent premature wear of the crimps (knuckles on woven wires). One of the pioneer products in this category was ONDAP GOMME made by the French manufacturer Giron. During the mid 90s, Major Wire Industries Ltd., a Quebec manufacturer, developed a “hybrid” self-cleaning screen cloth called Flex-Mat, without woven cross wires. In this product, the crimped longitudinal wires are held in place by polyurethane strips. Rather than locking (impeding) vibration over the support bars due to woven cross wires, polyurethane strips reduce vibration of longitudinal wires over the support bars, thus allowing vibration from hook to hook. Major Wire quickly started to promote this product as a high-performance screen that helped producers screen more in-specification material for less cost and not simply a problem solver. They claimed that the independent vibrating wires helped produce more product compared to a woven wire cloth with the same opening (aperture) and wire diameter. This higher throughput would be a direct result of the higher vibration frequency of each independent wire of the screen cloth (calculated in hertz) compared to the shaker vibration (calculated in RPM), accelerating the stratification of the material bed. Another benefit that helped the throughput increase is that hybrid self-cleaning screen media offered a better open area percentage than woven wire screen media. Due to its flat surface (no knuckles), hybrid self-cleaning screen media can use a smaller wire diameter for the same aperture than woven wire and still lasts as long, resulting in a greater opening percentage.
| Technology | Metallurgy | null |
12039054 | https://en.wikipedia.org/wiki/WHO%20Model%20List%20of%20Essential%20Medicines | WHO Model List of Essential Medicines | The WHO Model List of Essential Medicines (aka Essential Medicines List or EML), published by the World Health Organization (WHO), contains the medications considered to be most effective and safe to meet the most important needs in a health system. The list is frequently used by countries to help develop their own local lists of essential medicines. , more than 155 countries have created national lists of essential medicines based on the World Health Organization's model list. This includes both developed and developing countries.
The list is divided into core items and complementary items. The core items are deemed to be the most cost-effective options for key health problems and are usable with little additional health care resources. The complementary items either require additional infrastructure such as specially trained health care providers or diagnostic equipment or have a lower cost–benefit ratio. About 25% of items are in the complementary list. Some medications are listed as both core and complementary. While most medications on the list are available as generic products, being under patent does not preclude inclusion.
The first list was published in 1977 and included 208 medications. The WHO updates the list every two years. There are 306 medications in the 14th list in 2005, 410 in the 19th list in 2015, 433 in the 20th list in 2017, 460 in the 21st list in 2019, and 479 in the 22nd list in 2021. Various national lists contain between 334 and 580 medications. The Essential Medicines List (EML) was updated in July 2023 to its 23rd edition. This list contains 1200 recommendations for 591 drugs and 103 therapeutic equivalents.
A separate list for children up to 12 years of age, known as the WHO Model List of Essential Medicines for Children (EMLc), was created in 2007 and is in its 9th edition. It was created to make sure that the needs of children were systematically considered such as availability of proper formulations. Everything in the children's list is also included in the main list. The list and notes are based on the 19th to 23rd edition of the main list. Therapeutic alternatives with similar clinical performance are listed for some medicines and they may be considered for national essential medicines lists. The 9th Essential Medicines List for Children was updated in July 2023.
Note: An α indicates a medicine is on the complementary list.
Anaesthetics, preoperative medicines and medical gases
General anaesthetics and oxygen
Inhalational medicines
Halothane
Isoflurane
Nitrous oxide
Oxygen
Sevoflurane
Injectable medicines
Ketamine
Propofol
Thiopental
Local anaesthetics
Bupivacaine
Lidocaine
Lidocaine/epinephrine (lidocaine + epinephrine)
Complementary:
Ephedrine
Preoperative medication and sedation for short-term procedures
Atropine
Midazolam
Morphine
Medical gases
Oxygen
Medicines for pain and palliative care
Non-opioids and non-steroidal anti-inflammatory medicines (NSAIMs)
Acetylsalicylic acid (aspirin)
Ibuprofen
Paracetamol (acetaminophen)
Opioid analgesics
Codeine
Fentanyl
Morphine
Complementary:
Methadone
Medicines for other common symptoms in palliative care
Amitriptyline
Cyclizine
Dexamethasone
Diazepam
Docusate sodium
Fluoxetine
Haloperidol
Hyoscine butylbromide
Hyoscine hydrobromide
Lactulose
Loperamide
Metoclopramide
Midazolam
Ondansetron
Senna
Antiallergics and medicines used in anaphylaxis
Dexamethasone
Epinephrine (adrenaline)
Hydrocortisone
Loratadine
Prednisolone
Antidotes and other substances used in poisonings
Non-specific
Charcoal, activated
Specific
Acetylcysteine
Atropine
Calcium gluconate
Methylthioninium chloride (methylene blue)
Naloxone
Penicillamine
Prussian blue
Sodium nitrite
Sodium thiosulfate
Complementary:
Deferoxamine
Dimercaprol
Fomepizole
Sodium calcium edetate
Succimer
Medicines for diseases of the nervous system
Antiseizure medicines
Carbamazepine
Diazepam
Lamotrigine
Levetiracetam
Lorazepam
Magnesium sulfate
Midazolam
Phenobarbital
Phenytoin
Valproic acid (sodium valproate)
Complementary:
Ethosuximide
Levetiracetam
Valproic acid (sodium valproate)
Medicines for multiple sclerosis
Complementary:
Cladribine
Glatiramer acetate
Rituximab
Medicines for parkinsonism
Biperiden
Levodopa/carbidopa (levodopa + carbidopa)
Anti-infective medicines
Anthelminthics
Intestinal anthelminthics
Albendazole
Ivermectin
Levamisole
Mebendazole
Niclosamide
Praziquantel
Pyrantel
Antifilarials
Albendazole
Diethylcarbamazine
Ivermectin
Antischistosomals and other antinematode medicines
Praziquantel
Triclabendazole
Complementary:
Oxamniquine
Cysticidal medicines
Complementary:
Albendazole
Mebendazole
Praziquantel
Antibacterials
Access group antibiotics
Amikacin
Amoxicillin
Amoxicillin/clavulanic acid (amoxicillin + clavulanic acid)
Ampicillin
Benzathine benzylpenicillin
Benzylpenicillin
Cefalexin
Cefazolin
Chloramphenicol
Clindamycin
Cloxacillin
Doxycycline
Gentamicin
Metronidazole
Nitrofurantoin
Phenoxymethylpenicillin (penicillin V)
Procaine benzylpenicillin
Spectinomycin
Sulfamethoxazole/trimethoprim (sulfamethoxazole + trimethoprim)
Trimethoprim
Watch group antibiotics
Azithromycin
Cefixime
Cefotaxime
Ceftriaxone
Cefuroxime
Ciprofloxacin
Clarithromycin
Piperacillin/tazobactam (piperacillin + tazobactam)
Vancomycin
Complementary:
Ceftazidime
Meropenem
Vancomycin
Reserve group antibiotics
Reserve antibiotics are last-resort antibiotics. The EML antibiotic book was published in 2022.
Complementary:
Cefiderocol
Ceftazidime/avibactam (ceftazidime + avibactam)
Ceftolozane/tazobactam (ceftolozane + tazobactam)
Colistin
Fosfomycin
Linezolid
Meropenem/vaborbactam (meropenem + vaborbactam)
Plazomicin
Polymyxin B
Antileprosy medicines
Clofazimine
Dapsone
Rifampicin
Antituberculosis medicines
Ethambutol
Ethambutol/isoniazid/pyrazinamide/rifampicin (ethambutol + isoniazid + pyrazinamide + rifampicin)
Ethambutol/isoniazid/rifampicin (ethambutol + isoniazid + rifampicin)
Ethionamide
Isoniazid
Isoniazid/pyrazinamide/rifampicin (isoniazid + pyrazinamide + rifampicin)
Isoniazid/rifampicin (isoniazid + rifampicin)
Isoniazid/rifapentine (isoniazid + rifapentine)
Moxifloxacin
Pyrazinamide
Rifabutin
Rifampicin
Rifapentine
Complementary:
Amikacin
Amoxicillin/clavulanic acid (amoxicillin + clavulanic acid)
Bedaquiline
Clofazimine
Cycloserine
Delamanid
Ethionamide
Levofloxacin
Linezolid
Meropenem
Moxifloxacin
P-aminosalicylic acid (p-aminosalicylate sodium)
Pretomanid
Streptomycin
Antifungal medicines
Amphotericin B
Clotrimazole
Fluconazole
Flucytosine
Griseofulvin
Itraconazole
Nystatin
Voriconazole
Complementary:
Micafungin
Potassium iodide
Antiviral medicines
Antiherpes medicines
Aciclovir
Antiretrovirals
Nucleoside/nucleotide reverse transcriptase inhibitors
Abacavir
Lamivudine
Tenofovir disoproxil fumarate
Zidovudine
Non-nucleoside reverse transcriptase inhibitors
Efavirenz
Nevirapine
Protease inhibitors
Atazanavir/ritonavir (atazanavir + ritonavir)
Darunavir
Lopinavir/ritonavir (lopinavir + ritonavir)
Ritonavir
Integrase inhibitors
Dolutegravir
Raltegravir
Fixed-dose combinations of antiretroviral medicines
Abacavir/lamivudine (abacavir + lamivudine)
Dolutegravir/lamivudine/tenofovir (dolutegravir + lamivudine + tenofovir)
Efavirenz/emtricitabine/tenofovir
Efavirenz/lamivudine/tenofovir (efavirenz + lamivudine + tenofovir)
Emtricitabine/tenofovir (emtricitabine + tenofovir)
Lamivudine/zidovudine (lamivudine + zidovudine)
Medicines for prevention of HIV-related opportunistic infections
Isoniazid/pyridoxine/sulfamethoxazole/trimethoprim (isoniazid + pyridoxine + sulfamethoxazole + trimethoprim)
Other antivirals
Ribavirin
Valganciclovir
Complementary:
Oseltamivir
Valganciclovir
Antihepatitis medicines
Medicines for hepatitis B
Nucleoside/Nucleotide reverse transcriptase inhibitors
Entecavir
Tenofovir disoproxil fumarate
Medicines for hepatitis C
Pangenotypic direct-acting antiviral combinations
Daclatasvir
Daclatasvir/sofosbuvir (daclatasvir + sofosbuvir)
Glecaprevir/pibrentasvir (glecaprevir + pibrentasvir)
Ravidasvir
Sofosbuvir
Sofosbuvir/velpatasvir (sofosbuvir + velpatasvir)
Non-pangenotypic direct-acting antiviral combinations
Ledipasvir/sofosbuvir (ledipasvir + sofosbuvir)
Other antivirals for hepatitis C
Ribavirin
Antiprotozoal medicines
Antiamoebic and antigiardiasis medicines
Diloxanide
Metronidazole
Antileishmaniasis medicines
Amphotericin B
Meglumine antimoniate
Miltefosine
Paromomycin
Sodium stibogluconate
Antimalarial medicines
For curative treatment
Amodiaquine
Artemether
Artemether/lumefantrine (artemether + lumefantrine)
Artesunate
Artesunate/amodiaquine (artesunate + amodiaquine)
Artesunate/mefloquine (artesunate + mefloquine)
Artesunate/pyronaridine tetraphosphate (artesunate + pyronaridine tetraphosphate)
Chloroquine
Dihydroartemisinin/piperaquine phosphate (dihydroartemisinin + piperaquine phosphate)
Doxycycline
Mefloquine
Primaquine
Quinine
Sulfadoxine/pyrimethamine (sulfadoxine + pyrimethamine)
For chemoprevention
Amodiaquine + sulfadoxine/pyrimethamine (Co-packaged)
Chloroquine
Doxycycline
Mefloquine
Proguanil
Sulfadoxine/pyrimethamine (sulfadoxine + pyrimethamine)
Antipneumocystosis and antitoxoplasmosis medicines
Pyrimethamine
Sulfadiazine
Sulfamethoxazole/trimethoprim (sulfamethoxazole + trimethoprim)
Complementary:
Pentamidine
Antitrypanosomal medicines
African trypanosomiasis
Fexinidazole
Medicines for the treatment of 1st stage African trypanosomiasis
Pentamidine
Suramin sodium
Medicines for the treatment of 2nd stage African trypanosomiasis
Eflornithine
Melarsoprol
Nifurtimox
Complementary:
Melarsoprol
American trypanosomiasis
Benznidazole
Nifurtimox
Medicines for ectoparasitic infections
Ivermectin
Medicines for Ebola virus disease
Ansuvimab
Atoltivimab/maftivimab/odesivimab (atoltivimab + maftivimab + odesivimab)
Medicines for COVID-19
No listings in this section.
Antimigraine medicines
For treatment of acute attack
Acetylsalicylic acid (aspirin)
Ibuprofen
Paracetamol (acetaminophen)
Sumatriptan
For prophylaxis
Propranolol
Immunomodulators and antineoplastics
Immunomodulators for non-malignant disease
Complementary:
Adalimumab
Azathioprine
Ciclosporin
Tacrolimus
Antineoplastics and supportive medicines
Cytotoxic medicines
Complementary:
Arsenic trioxide
Asparaginase
Bendamustine
Bleomycin
Calcium folinate (leucovorin calcium)
Capecitabine
Carboplatin
Chlorambucil
Cisplatin
Cyclophosphamide
Cytarabine
Dacarbazine
Dactinomycin
Daunorubicin
Docetaxel
Doxorubicin
Doxorubicin (as pegylated liposomal)
Etoposide
Fludarabine
Fluorouracil
Gemcitabine
Hydroxycarbamide (hydroxyurea)
Ifosfamide
Irinotecan
Melphalan
Mercaptopurine
Methotrexate
Oxaliplatin
Paclitaxel
Pegaspargase
Procarbazine
Realgar Indigo naturalis formulation
Tioguanine
Vinblastine
Vincristine
Vinorelbine
Targeted therapies
Complementary:
All-trans retinoic acid (tretinoin) (ATRA)
Bortezomib
Dasatinib
Erlotinib
Everolimus
Ibrutinib
Imatinib
Nilotinib
Rituximab
Trastuzumab
Immunomodulators
Complementary:
Filgrastim
Lenalidomide
Nivolumab
Pegfilgrastim
Thalidomide
Hormones and antihormones
Complementary:
Abiraterone
Anastrozole
Bicalutamide
Dexamethasone
Hydrocortisone
Leuprorelin
Methylprednisolone
Prednisolone
Tamoxifen
Supportive medicines
Complementary:
Allopurinol
Mesna
Rasburicase
Zoledronic acid
Therapeutic foods
Ready-to-use therapeutic food
Medicines affecting the blood
Antianaemia medicines
Ferrous salt
Ferrous salt/folic acid (ferrous salt + folic acid)
Folic acid
Hydroxocobalamin
Complementary:
Erythropoiesis-stimulating agents
Medicines affecting coagulation
Dabigatran
Enoxaparin
Heparin sodium
Phytomenadione
Protamine sulfate
Tranexamic acid
Warfarin
Complementary:
Desmopressin
Heparin sodium
Protamine sulfate
Warfarin
Other medicines for haemoglobinopathies
Deferasirox
Complementary:
Deferoxamine
Hydroxycarbamide (hydroxyurea)
Blood products of human origin and plasma substitutes
Blood and blood components
Cryoprecipitate, pathogen-reduced
Fresh frozen plasma
Platelets
Red blood cells
Whole blood
Plasma-derived medicines
Human immunoglobulins
Rho(D) immune globulin (anti-D immunoglobulin)
Anti-rabies immunoglobulin
Anti-tetanus immunoglobulin
Complementary:
Normal immunoglobulin
Blood coagulation factors
Complementary:
Coagulation factor VIII
Coagulation factor IX
Plasma substitutes
Dextran 70
Cardiovascular medicines
Antianginal medicines
Bisoprolol
Glyceryl trinitrate
Isosorbide dinitrate
Verapamil
Antiarrhythmic medicines
Bisoprolol
Digoxin
Epinephrine (adrenaline)
Lidocaine
Verapamil
Complementary:
Amiodarone
Antihypertensive medicines
Amlodipine
Bisoprolol
Enalapril
Hydralazine
Hydrochlorothiazide
Lisinopril/amlodipine (lisinopril + amlodipine)
Lisinopril/hydrochlorothiazide (lisinopril + hydrochlorothiazide)
Losartan
Methyldopa
Telmisartan/amlodipine (telmisartan + amlodipine)
Telmisartan/hydrochlorothiazide (telmisartan + hydrochlorothiazide)
Complementary:
Sodium nitroprusside
Medicines used in heart failure
Bisoprolol
Digoxin
Enalapril
Furosemide
Hydrochlorothiazide
Losartan
Spironolactone
Complementary:
Digoxin
Dopamine
Antithrombotic medicines
Anti-platelet medicines
Acetylsalicylic acid (aspirin)
Clopidogrel
Thrombolytic medicines
Complementary:
Alteplase
Streptokinase
Lipid-lowering agents
Simvastatin
Fixed-dose combinations for prevention of atherosclerotic cardiovascular disease
Acetylsalicylic acid/atorvastatin/ramipril (acetylsalicylic acid + atorvastatin + ramipril)
Acetylsalicylic acid/simvastatin/ramipril/atenolol/hydrochlorothiazide (acetylsalicylic acid + simvastatin + ramipril + atenolol + hydrochlorothiazide)
Atorvastatin/perindopril/amlodipine (atorvastatin + perindopril + amlodipine)
Dermatological medicines (topical)
Antifungal medicines
Miconazole
Selenium sulfide
Sodium thiosulfate
Terbinafine
Anti-infective medicines
Mupirocin
Potassium permanganate
Silver sulfadiazine
Anti-inflammatory and antipruritic medicines
Betamethasone
Calamine
Hydrocortisone
Medicines affecting skin differentiation and proliferation
Benzoyl peroxide
Calcipotriol
Coal tar
Fluorouracil
Podophyllum resin
Salicylic acid
Urea
Complementary:
Methotrexate
Scabicides and pediculicides
Benzyl benzoate
Permethrin
Diagnostic agents
Ophthalmic medicines
Fluorescein
Tropicamide
Radiocontrast media
Amidotrizoate
Barium sulfate
Iohexol
Complementary:
Barium sulfate
Meglumine iotroxate
Antiseptics and disinfectants
Antiseptics
Chlorhexidine
Ethanol
Povidone iodine
Disinfectants
Alcohol based hand rub
Chlorine base compound
Chloroxylenol
Glutaral
Diuretics
Amiloride
Furosemide
Hydrochlorothiazide
Mannitol
Spironolactone
Complementary:
Hydrochlorothiazide
Mannitol
Spironolactone
Gastrointestinal medicines
Complementary:
Pancreatic enzymes
Antiulcer medicines
Omeprazole
Ranitidine
Antiemetic medicines
Dexamethasone
Metoclopramide
Ondansetron
Complementary:
Aprepitant
Anti-inflammatory medicines
Sulfasalazine
Complementary:
Hydrocortisone
Prednisolone
Laxatives
Senna
Medicines used in diarrhoea
Oral rehydration salts + zinc sulfate (Co-packaged)
Oral rehydration
Oral rehydration salts
Medicines for diarrhoea
Zinc sulfate
Medicines for endocrine disorders
Adrenal hormones and synthetic substitutes
Fludrocortisone
Hydrocortisone
Androgens
Complementary:
Testosterone
Estrogens
No listings in this section.
Progestogens
Medroxyprogesterone acetate
Medicines for diabetes
Insulins
Insulin injection (soluble)
Intermediate-acting insulin
Long-acting insulin analogues
Oral hypoglycaemic agents
Empagliflozin
Gliclazide
Metformin
Complementary:
Metformin
Medicines for hypoglycaemia
Glucagon
Complementary:
Diazoxide
Thyroid hormones and antithyroid medicines
Levothyroxine
Potassium iodide
Methimazole
Propylthiouracil
Complementary:
Lugol's solution
Methimazole
Potassium iodide
Propylthiouracil
Medicines for disorders of the pituitary hormone system
Cabergoline
Complementary:
Octreotide
Immunologicals
Diagnostic agents
Tuberculin, purified protein derivative (PPD)
Sera, immunoglobulins and monoclonal antibodies
Anti-rabies virus monoclonal antibodies
Antivenom immunoglobulin
Diphtheria antitoxin
Equine rabies immunoglobulin
Vaccines
Recommendations for all
BCG vaccine
Diphtheria vaccine
Haemophilus influenzae type b vaccine
Hepatitis B vaccine
Human papilloma virus (HPV) vaccine
Measles vaccine
Pertussis vaccine
Pneumococcal vaccine
Poliomyelitis vaccine
Rotavirus vaccine
Rubella vaccine
Tetanus vaccine
Recommendations for certain regions
Japanese encephalitis vaccine
Tick-borne encephalitis vaccine
Yellow fever vaccine
Recommendations for some high-risk populations
Cholera vaccine
Dengue vaccine
Hepatitis A vaccine
Meningococcal meningitis vaccine
Rabies vaccine
Typhoid vaccine
Recommendations for immunization programmes with certain characteristics
Influenza vaccine (seasonal)
Mumps vaccine
Varicella vaccine
Muscle relaxants (peripherally-acting) and cholinesterase inhibitors
Atracurium
Neostigmine
Suxamethonium
Vecuronium
Complementary:
Pyridostigmine
Vecuronium
Ophthalmological preparations
Anti-infective agents
Aciclovir
Azithromycin
Erythromycin
Gentamicin
Natamycin
Ofloxacin
Tetracycline
Anti-inflammatory agents
Prednisolone
Local anesthetics
Tetracaine
Miotics and antiglaucoma medicines
Acetazolamide
Latanoprost
Pilocarpine
Timolol
Mydriatics
Atropine
Complementary:
Epinephrine (adrenaline)
Anti-vascular endothelial growth factor (VEGF) preparations
Complementary:
Bevacizumab
Medicines for reproductive health and perinatal care
Contraceptives
Oral hormonal contraceptives
Ethinylestradiol/levonorgestrel (ethinylestradiol + levonorgestrel)
Ethinylestradiol/norethisterone (ethinylestradiol + norethisterone)
Levonorgestrel
Ulipristal
Injectable hormonal contraceptives
Estradiol cypionate/medroxyprogesterone acetate (estradiol cypionate + medroxyprogesterone acetate)
Medroxyprogesterone acetate
Norethisterone enantate
Intrauterine devices
Copper-containing device
Levonorgestrel-releasing intrauterine system
Barrier methods
Condoms
Diaphragms
Implantable contraceptives
Etonogestrel-releasing implant
Levonorgestrel-releasing implant
Intravaginal contraceptives
Ethinylestradiol/etonogestrel (ethinylestradiol + etonogestrel)
Progesterone vaginal ring
Ovulation inducers
Complementary:
Clomifene
Letrozole
Uterotonics
Carbetocin
Ergometrine
Mifepristone + misoprostol (Co-packaged)
Misoprostol
Oxytocin
Antioxytocics (tocolytics)
Nifedipine
Other medicines administered to the mother
Dexamethasone
Multiple micronutrient supplement
Tranexamic acid
Medicines administered to the neonate
Caffeine citrate
Chlorhexidine
Complementary:
Ibuprofen
Prostaglandin E1
Surfactant
Peritoneal dialysis solution
Complementary:
Intraperitoneal dialysis solution (of appropriate composition)
Medicines for mental and behavioural disorders
Medicines used in psychotic disorders
Fluphenazine
Haloperidol
Olanzapine
Paliperidone
Risperidone
Complementary:
Clozapine
Medicines used in mood disorders
Medicines used in depressive disorders
Amitriptyline
Fluoxetine
Medicines used in bipolar disorders
Carbamazepine
Lithium carbonate
Quetiapine
Valproic acid (sodium valproate)
Medicines for anxiety disorders
Diazepam
Fluoxetine
Medicines used for obsessive compulsive disorders
Clomipramine
Fluoxetine
Medicines for disorders due to psychoactive substance use
Medicines for alcohol use disorders
Acamprosate calcium
Naltrexone
Medicines for nicotine use disorders
Bupropion
Nicotine replacement therapy (NRT)
Varenicline
Complementary:
Methadone
Medicines acting on the respiratory tract
Antiasthmatic medicines and medicines for chronic obstructive pulmonary disease
Budesonide
Budesonide/formoterol (budesonide + formoterol)
Epinephrine (adrenaline)
Ipratropium bromide
Salbutamol
Tiotropium
Solutions correcting water, electrolyte and acid-base disturbances
Oral
Oral rehydration salts
Potassium chloride
Parenteral
Glucose
Glucose with sodium chloride
Potassium chloride
Sodium chloride
Sodium hydrogen carbonate
Sodium lactate, compound solution (Ringer's lactate solution)
Miscellaneous
Water for injection
Vitamins and minerals
Ascorbic acid
Calcium
Colecalciferol
Ergocalciferol
Iodine
Multiple micronutrient powder
Nicotinamide
Pyridoxine
Retinol
Riboflavin
Thiamine
Complementary:
Calcium gluconate
Ear, nose and throat medicines
Acetic acid
Budesonide
Ciprofloxacin
Xylometazoline
Medicines for diseases of joints
Medicines used to treat gout
Allopurinol
Disease-modifying anti-rheumatic drugs (DMARDs)
Chloroquine
Complementary:
Azathioprine
Hydroxychloroquine
Methotrexate
Penicillamine
Sulfasalazine
Medicines for juvenile joint diseases
Complementary:
Acetylsalicylic acid (aspirin)
Adalimumab
Methotrexate
Triamcinolone hexacetonide
Dental medicines and preparations
Fluoride
Glass ionomer cement
Resin-based composite (low-viscosity)
Resin-based composite (high-viscosity)
Silver diamine fluoride
| Biology and health sciences | General concepts_2 | Health |
7391204 | https://en.wikipedia.org/wiki/Transport%20hub | Transport hub | A transport hub is a place where passengers and cargo are exchanged between vehicles and/or between transport modes. Public transport hubs include railway stations, rapid transit stations, bus stops, tram stops, airports, and ferry slips. Freight hubs include classification yards, airports, seaports, and truck terminals, or combinations of these. For private transport by car, the parking lot functions as an unimodal hub.
History
Historically, an interchange service in the scheduled passenger air transport industry involved a "through plane" flight operated by two or more airlines where a single aircraft was used with the individual airlines operating it with their own flight crews on their respective portions of a direct, no-change-of-plane multi-stop flight. In the U.S., a number of air carriers including Alaska Airlines, American Airlines, Braniff International Airways, Continental Airlines, Delta Air Lines, Eastern Airlines, Frontier Airlines (1950-1986), Hughes Airwest, National Airlines (1934-1980), Pan Am, Trans World Airlines (TWA), United Airlines and Western Airlines previously operated such cooperative "through plane" interchange flights on both domestic and/or international services with these schedules appearing in their respective system timetables.
Delta Air Lines pioneered the hub and spoke system for aviation in 1955 from its hub in Atlanta, Georgia, United States, in an effort to compete with Eastern Air Lines. FedEx adopted the hub and spoke model for overnight package delivery during the 1970s. When the United States airline industry was deregulated in 1978, Delta's hub and spoke paradigm was adopted by several airlines. Many airlines around the world operate hub-and-spoke systems facilitating passenger connections between their respective flights.
Public transport
Intermodal passenger transport hubs in public transport include bus stations, railway stations and metro stations, while a major transport hub, often
multimodal (bus and rail), may be referred to as a transport centre or, in American English, as a transit center. Sections of city streets that are devoted to functioning as transit hubs are referred to as transit malls. In cities with a central station, that station often also functions as a transport hub in addition to being a railway station.
Journey planning involving transport hubs is more complicated than direct trips, as journeys will typically require a transfer at the hub. Modern electronic journey planners for public transport have a digital representation of both the stops and transport hubs in a network, to allow them to calculate journeys that include transfers at hubs.
Airports
Airports have a twofold hub function. First, they concentrate passenger traffic into one place for onward transportation. This makes it important for airports to be connected to the surrounding transport infrastructure, including roads, bus services, and railway and rapid transit systems. Secondly some airports function as intra-modular hubs for the airlines, or airline hubs. This is a common strategy among network airlines who fly only from limited number of airports and usually will make their customers change planes at one of their hubs if they want to get between two cities the airline does not fly directly between.
Airlines have extended the hub-and-spoke model in various ways. One method is to create additional hubs on a regional basis, and to create major routes between the hubs. This reduces the need to travel long distances between nodes that are close together. Another method is to use focus cities to implement point-to-point service for high traffic routes, bypassing the hub entirely.
Freight
There are usually three kinds of freight hubs: sea-road, sea-rail, and road-rail, though they can also be sea-road-rail. With the growth of containerization, intermodal freight transport has become more efficient, often making multiple legs cheaper than through services—increasing the use of hubs.
| Technology | Concepts of ground transport | null |
351249 | https://en.wikipedia.org/wiki/Apiary | Apiary | An apiary (also known as a bee yard) is a location where beehives of honey bees are kept. Apiaries come in many sizes and can be rural or urban depending on the honey production operation. Furthermore, an apiary may refer to a hobbyist's hives or those used for commercial or educational usage. It can also be a wall-less, roofed structure, similar to a gazebo which houses hives, or an enclosed structure with an opening that directs the flight path of the bees.
History
Apiaries have been found in ancient Egypt from prior to 2422 BCE where hives were constructed from moulded mud. Throughout history apiaries and bees have been kept for honey and pollination purposes all across the globe. Due to the definition of apiary as a location where hives are kept, its history can be traced as far back as that of beekeeping itself.
Etymology
The first known usage of the word "apiary" was in 1654. The base of the word comes from the Latin word "apis" meaning "bee", leading to "apiarium" or "beehouse" and eventually "apiary." Beekeepers may be referred to as "apiarists" or "ones who tend apiaries."
Structure
Apiaries may vary by location and according to the needs of the individual operation. Typically, apiaries are composed of several individual hives. For more information on specific hive structures see the beekeeping and beehive articles. In the case of urban beekeeping, hives are often located on high ground, which requires less space than hives located at lesser altitudes. To direct the bees' path of flight in populous urban areas, beekeepers often construct tall fences which force the bees to fly higher and widen their search for food or place the hives in an enclosed apiary with an opening that directs bees' flight path up overhead.
Location
Apiaries are usually situated on high ground in order to avoid moisture collection, though in proximity to a consistent water source—whether natural or man-made—to ensure the bees' access. Additionally, ample nectar supplies for the bees as well as relatively large amounts of sun are considered. They are often situated close to orchards, farms, and public gardens, which require frequent pollination to develop a positive feedback loop between the bees and their food sources. This also economizes on the bees' pollination and the plants' supply of nectar.
An apiary may have hive management objectives other than honey production, including queen rearing and mating. In the northern hemisphere, east and south facing locations with full morning sun are preferred. In hot climates, shade is needed and may have to be artificially provided if trees are not present. Other factors include air and water drainage and accessibility by truck, distance from phobic people, and protection from vandalism.
In the USA there are beekeepers—from hobbyists to commercial—in every state. The most lucrative areas for American honey production are Florida, Texas, California, and the Upper Midwest. For paid pollination, the main areas are California, the Pacific Northwest, the Great Lakes States, and the Northeast. Rules and regulations by local ordinances and zoning laws also affect apiaries.
In recent years US honey production has dropped and the U.S. imports 16% of the world's honey. Internationally, the largest honey producing exporters are China, Germany, and Mexico. As in the United States the location of apiaries varies internationally depending on available resources and the operational need. For more information on nation-specific beekeeping see their respective articles, such as the Beekeeping in Nepal article.
Size
Apiary size refers not only to the spatial size of the apiary, but also to the number of bee families and bees by weight. With ample space there is no limit to the number of hives or bee families which can be housed in an apiary. The larger the number of hives held in an apiary the higher the yield of honey relative to resources, often resulting in apiaries growing with time and experience. Additionally a higher number of hives within an apiary can increase the quality of the honey produced. Depending on the nectar and pollen sources in a given area, the maximum number of hives that can be placed in one apiary can vary. If too many hives are placed into an apiary, the hives compete with each other for scarce resources. This can lead to lower honey, flower pollen and bee bread yields, as well as higher transmission of disease and robbing.
The size of an apiary is determined by not only the resources available but also by the variety of honey being cultivated, with more complex types generally cultivated in smaller productions. For more specific details on varieties see the classification portion of the honey article. The purpose of the apiary also affects size: apiaries are kept by commercial and local honey producers, as well as by universities, research facilities, and local organizations. Many such organizations provide community programming and educational opportunities. This results in varying sizes of apiaries depending on usage characteristics.
The maximum size of a permanent apiary or bee yard may depend on the type of bee as well. Some honey bee species fly farther than others. A circle around an apiary with a three-mile (5 km) foraging radius covers 28 square miles (73 km2). A good rule of thumb is to have no more than 25–35 hives in a permanent apiary, although migrating beekeepers may temporarily place one hundred hives into a location with a good nectar flow.
Disease and decline
Apiaries may decline due to a scarcity of resources which can lead to robbing of nearby hives. This is especially an issue in urban areas where there may be a limited amount of resources for bees and a large number of hives may be affected.
Apiaries may suffer from a wide variety of diseases and infestations. Throughout history apiaries and bees have been kept for honey and pollination purposes all across the globe. Due to the definition of apiary as a location where hives are kept its history can be traced as far back as that of beekeeping itself. In recent years Colony Collapse Disorder due to pesticide resistant mites have ravaged bee populations. Beyond mites there are a wide variety of diseases which may affect the hives and lead to the decline or collapse of a colony. For this reason many beekeepers choose to keep apiaries of limited size to avoid mass infection or infestation. For more information on diseases which affect bee populations see the list of diseases of the honey bee.
| Technology | Buildings and infrastructure | null |
351263 | https://en.wikipedia.org/wiki/Indian%20summer | Indian summer | An Indian summer is a period of unseasonably warm, dry weather that sometimes occurs in autumn in temperate regions of the northern hemisphere. Several sources describe a true Indian summer as not occurring until after the first frost, or more specifically the first "killing frost".
Etymology
The late 19th-century lexicographer Albert Matthews made an exhaustive search of early American literature in an attempt to discover who coined the expression. The earliest reference he found dated to 1851. He also found the phrase in a letter written in England in 1778, but discounted that as a coincidental use of the phrase.
Later research showed that the earliest known reference to Indian summer in its current sense occurs in an essay written in the United States around 1778 by J. Hector St. John de Crevecœur, describing the character of autumn and implying the common usage of the expression
The essay was first published in French around 1788, but remained unavailable in the United States until the 1920s.
Although the exact origins of the term are uncertain, it was perhaps so-called because it was first noted in regions inhabited by Native Americans, or because the natives first described it to Europeans, or it had been based on the warm and hazy conditions in autumn when Native Americans hunted. John James Audubon wrote about "The Indian Summer that extraordinary phenomenon of North America" in his journal on November 20, 1820. He mentions the "constant Smoky atmosphere" and how the smoke irritates his eyes. Audubon suspects that the condition of the air was caused by "Indians, firing the Prairies of the West". Audubon also mentions in many other places in his writings the reliance Native Americans had on fire. At no point does Audubon relate an Indian summer to warm temperatures during the cold seasons.
Because the warm weather is not a permanent gift, a connection has been made to the pejorative term Indian giver. Native-American legends mention the god or "Life-Giver" bestowing warm autumnal weather to various warriors or peoples, enabling them to survive after great misfortune, such as loss of crops.
Usage
Weather historian William R. Deedler wrote that "Indian summer" can be defined as "any spell of warm, quiet, hazy weather that may occur in October or November", though he noted that he "was surprised to read that Indian summers have been given credit for warm spells as late as December and January". Deedler also noted that some writers use Indian summer in reference to the weather in only New England, "while others have stated it happens over most of the United States, even along the Pacific coast".
In literature and history, the term is sometimes used metaphorically. The title of Van Wyck Brooks' New England: Indian Summer (1940) suggests an era of inconsistency, infertility, and depleted capabilities, a period of seemingly robust strength that is only an imitation of an earlier season of actual strength. William Dean Howells' 1886 novel Indian Summer uses the term to mean a time when one may recover some of the happiness of youth. The main character, jilted as a young man, leads a solitary life until he rediscovers romance in early middle age.
In British English, the term is used in the same way as in North America. In the UK it was first used in the early 19th century, there was an early example in The Guardian which ran an article explaining the phrase "Indian summer" to its readers in 1837, written by someone who had lived in the US but questioned whether Native Americans had influenced the origins. The UK Met Office Meteorological Glossary published in 1916 defines an Indian summer "a warm, calm spell of weather occurring in autumn, especially in October and November", while The Indian Summer of a Forsyte is the metaphorical title of the 1918 second volume of The Forsyte Saga by John Galsworthy. However, early 20th-century climatologists Gordon Manley and Hubert Lamb used it only when referring to the American phenomenon, and the expression did not gain wide currency in Great Britain until the 1950s. In former times, variations of "Saint Martin's summer" were widely used across Europe to describe warm weather surrounding autumn feast days of St. Martin and Saint Luke.
In the English translation of Boris Pasternak's Doctor Zhivago, the term is used to describe the unseasonably warm weather leading up to the October Revolution.
Other names and similar phenomena
Similar weather conditions with local variations also exist. A warm period in autumn is called ("old women's summer") in Germany, Austria, Switzerland, Lithuania, Hungary (), Estonia (), and in a number of Slavic-language countries—for example, in the Czech Republic, Ukraine, Poland, Slovakia, Russia and Slovenia, – it is known as "(old) women's summer" (; ; , , ; ; ). In Bulgaria, it is known as "gypsy summer" or "poor man's summer", and in Serbia it is known as "Miholjsko leto" because Saint Michael or "Miholjdan" is celebrated on October 12. In Sweden, there's "Brittsommar" (out of "Birgitta" and "Britta", having their name days around the time, on October 7) and/or "Indiansommar" as a direct translation from English. In Finland, the period is today called , a direct translation, but historically a warm period in autumn was named after Bartholomew ( or Perttu), his saint day being in late August. In Irish, the phenomenon is called ("little autumn of the geese"). In Spain is also known as ("little summer of the quince tree") or ("the sun of the quince tree").
In temperate parts of South America—such as southernmost Brazil, Argentina, Chile and Uruguay—the phenomenon is known as "Veranico", "Veranito" or "Veranillo" (literally, "little summer"), and usually occurs in early autumn, between late April and mid-May, when it is known as "Veranico de Maio" ("May's little summer") or as "Veranito de San Juan" ("Saint John's little summer"). Its onset and duration are directly associated with the occurrence of El Niño.
In other countries, it is associated with autumnal name days or saint days, such as Teresa of Ávila (Portugal, Spain and France), St. Martin's Summer (Spain, France, Italy, Portugal and Malta), St. Michael's summer ( in Spain, , Serbia, Montenegro and Bosnia and Herzegovina), St. Martin's Day (Netherlands and Italy), St. Demetrius (Greece and Cyprus), Bridget of Sweden in Sweden, and Saint Michael the Archangel in Wales (). In Turkey, it is called , meaning "pastrami summer", since the month of November was considered to be the best time to make (the meat that, though slightly different, pastrami originated from).
The American Meteorological Society (AMS) also notes that a similar phenomenon may be referred to poetically as halcyon days, a term that originated in Greek mythology. Halcyon days in Greece take place in winter, usually 16–31 of January and last around 4–7 days with extremely warm and sunny days. "All-hallown summer" or "All Saints' summer" is also referenced in English folklore and by Shakespeare, but its use appears to have died out.
In media
Board games
Indian Summer, designed by Uwe Rosenberg, is named and themed after the event, and involves players placing leaf-filled tiles on the forest floor.
Books
Engine Summer written by John Crowley in 1979, is named after and refers to the event, with the spelling changed to reflect the post-apocalyptic setting of the book.
Indian Summer by John Knowles, published in 1966.
Indian Summer was written by Adalbert Stifter in 1857.
Indian Summer was written by William Dean Howells in 1886.
Indian Summer: The Secret History of the End of an Empire was written by Alex von Tunzelmann in 2007.
Indian Summer: The Tragic Story of Louis Francis Sockalexis, the First Native American in Major League Baseball was written by Brian McDonald in 2003.
The graphic novel Indian Summer was written by Hugo Pratt and illustrated by Milo Manara in 1983.
The Indian Summer Of English Chivalry written by Arthur Ferguson in 1960.
Comics
Indian Summer, Hugo Pratt, Nantier Beall Minoustchine, October 1, 1993.
Injun Summer, John T. McCutcheon, Chicago Tribune, September 30, 1907.
Music
The Victor Herbert piano solo with this title dates to 1919. It received an Al Dubin lyric in the 1930s and was recorded by several pop singers and dance bands.
In 1945, Coleman Hawkins recorded a jazz version of the Victor Herbert/Al Dubin tune on tenor sax.
In 1966, The Doors recorded their original song "Indian Summer" (Morrison/Krieger), which was released on their 1970 album Morrison Hotel.
In 1969, Brewer & Shipley recorded their own song "Indian Summer", for the Weeds album.
In 1975, Joe Dassin recorded the song "Indian Summer" in French, English, Spanish and German. "L'Été indien" was based on the song "Africa" by Toto Cutugno, hence the subtitle "L'Été indien (Africa)" on some single releases. It went on to become Dassin's biggest hit, selling almost 2 million copies worldwide. Nancy Sinatra and Lee Hazlewood released an English language cover of the song as a single in 1976.
In 1977, Poco released the album Indian Summer, which contained the title track written by Paul Cotton.
Jay Ferguson's 1977 song "Thunder Island" contains the passage "She was the color of the Indian Summer".
In 1978, Joe Walsh recorded his song "Indian Summer" for the album But Seriously, Folks....
In 1981, Al Stewart released his song "Indian Summer" on his first live album Live/Indian Summer.
In 1983, Belle Stars released a single called "Indian Summer". It also features on the Belle Stars album.
In 1983, Per Gessle released an instrumental song called "Indiansommar" (Swedish for Indian summer) on his self-titled debut album.
In 1984, U2 included "Indian Summer Sky" on their The Unforgettable Fire album.
In 1985, Larry Gatlin and Barry Gibb wrote their song "Indian Summer", which was released on the Larry Gatlin & The Gatlin Brothers album Smile (1985), as performed by Larry, Barry and Roy Orbison.
In 1985, Channel 3 (band) included "Indian Summer" on their Last Time I Drank album.
In 1987, the band Opal released their version of The Doors song on the Chemical Imbalance Limited Edition 45 (#003).
In 1987, The Dream Academy recorded their song "Indian Summer" for the album Remembrance Days.
In 1988, Beat Happening released the Calvin Johnson penned "Indian Summer" on their album Jamboree.
In 1992, Go West released an album called Indian Summer.
In 1992, the Victor Herbert/Al Dubin tune was recorded by Tony Bennett for his Frank Sinatra tribute album, Perfectly Frank.
In 1992, The Rippingtons released "Indian Summer" as the fourth track on their album Weekend in Monaco.
In 1993, Paul Westerberg released "First Glimmer," a song that references an Indian Summer.
In 1993, Luna released their version of the Beat Happening song on their EP Slide.
In 1993, the emo band Indian Summer was formed in Oakland, California. They disbanded in 1994.
In 2002, Pedro the Lion released the David Bazan penned "Indian Summer" on their album Control.
In 2004, Carbon Leaf released a collection of all-original songs on their album, Indian Summer released on Vanguard Records.
In 2004, Tori Amos recorded her song "Indian Summer" for the EP Scarlet's Hidden Treasures.
In 2007, Ben Gibbard's version of the "Beat Happening" song was included on the Kurt Cobain About a Son: Music from the Motion Picture soundtrack.
In 2007, classical composer Pyarelal Sharma wrote Indian Summer: 8 Enchanting Pieces for String Quartet.
In 2007, jazz musician Dave Brubeck released his first solo piano album in 50 years on Telark, called "Indian Summer", after his version of the title song by Victor Herbert and Al Dubin.
In 2007, Manic Street Preachers released their song "Indian Summer" as the third single released from their album Send Away the Tigers.
In 2009, country duo Brooks and Dunn released their own "Indian Summer", as the lead single to their fifth greatest hits package, #1s… and Then Some.
In 2009, Mandy Moore released the album Amanda Leigh which includes the song "Indian Summer" that she co-wrote with Mike Viola and Inara George.
In 2010, Australian record producer Gabriel Gleeson began releasing electronic music and performing under the name Indian Summer.
In 2011, Loaded (sometimes called Duff McKagan's Loaded) released their song "Indian Summer" on the album called The Taking.
In 2013, Stereophonics released the Kelly Jones penned "Indian Summer", as the second single from their album Graffiti on the Train.
In 2014, Tyler Hilton released the album Indian Summer, containing his self-penned title track.
In 2015, Jai Wolf released his debut single "Indian Summer" on the Foreign Family Collective label.
In 2018, Dutch singer Sharon den Adel released the song Indian Summer under her My Indigo project.
In 2020, No Germ Candy released their version of the Beat Happening as a b-side to the Straight Talk single.
Katy Perry's 2009 song "Thinking of You" contains the passage "You're like an Indian summer in the middle of a winter".
Other jazz versions based on the Victor Herbert tune with Al Dubin lyrics were recorded by the Ginny Simms with Kay Kyser & his Orchestra (December 1939 recording for Columbia 78rpm single), Gene Krupa Orchestra (recorded live on radio, January 1940), Bing Crosby (on Bing Crosby – Victor Herbert 7" 45rpm box set, for Decca in 1950), Lee Konitz and Billy Bauer (recorded for Prestige on 1951 Lee Konitz: The New Sounds 10" and 1956 Conception LP), Stan Getz on the Stan Getz Quartets LP (recorded June 1949 for Prestige LP in 1955), The Hi-Lo's from their On Hand LP (Starlite 1956), Joe Puma with Bill Evans (from the album Joe Puma Jazz Trio and Quartet, on Jubilee, 1957), Dave Brubeck on his first solo piano album Dave Brubeck Plays and Plays and... (Fantasy Records 1957), Ella Fitzgerald with the Count Basie Orchestra featuring the Tommy Flanagan Trio (recorded live in 1972), Paul Desmond (1973 on Skylark), and Sarah Vaughan with the Count Basie Orchestra (on Send in the Clowns 1974).
Sidney Bechet recorded a jazz version of the Victor Herbert/Al Dubin tune on soprano sax in 1940.
The Glenn Miller Big Band Orchestra version of Victor Herbert and Al Dubin's tune with vocalist Ray Eberle, rose to number 8 from late 1939 into 1940.
The Victor Herbert/Al Dubin tune was a number 1 hit for Tommy Dorsey's Big Band Orchestra with Jack Leonard on vocals in 1939.
The Victor Herbert/Al Dubin tune was recorded by Frank Sinatra on his album with Duke Ellington, Francis A, and Edward K., in 1968.
Victor Herbert composed the song "Indian Summer" in 1919 for classical orchestra and Al Dubin wrote lyrics in 1939.
Painting
In 1875, Józef Chełmoński painted a picture Indian Summer with a wide landscape panorama.
In 1922, Willard Leroy Metcalf painted Indian Summer, Vermont
Poetry
William Wilfred Campbell's poem "Indian Summer".
Emily Dickinson wrote some 20 poems about Indian Summer, including "These Are the Days When Birds Come Back", "The Gentian Weaves Her Fringes" and "There Is a June When Corn Is Cut".
Kate Harrington wrote "Legend of the Indian Summer".
Oliver Wendell Holmes wrote "Our Indian Summer".
Vachel Lindsay wrote "An Indian Summer Day on the Prairie".
Henry Wadsworth Longfellow's poem The Song of Hiawatha (1855) mentions "the tender Indian Summer"
Jayanta Mahapatra wrote "Indian Summer".
Barry Middleton wrote "Indian Summer".
Dorothy Parker wrote her own "Indian Summer".
Robert William Service wrote "My Indian Summer".
Lydia Sigourney's poem "The Indian Summer" was published in the volume Illustrated Poems, 1849.
Alma Luz Villanueva wrote "Indian Summer Ritual".
| Physical sciences | Meteorology: General | Earth science |
351423 | https://en.wikipedia.org/wiki/Foehn%20wind | Foehn wind | A Foehn, or Föhn (, , ), is a type of dry, relatively warm downslope wind in the lee of a mountain range. It is a rain shadow wind that results from the subsequent adiabatic warming of air that has dropped most of its moisture on windward slopes (see orographic lift). As a consequence of the different adiabatic lapse rates of moist and dry air, the air on the leeward slopes becomes warmer than equivalent elevations on the windward slopes.
Foehn winds can raise temperatures by as much as in just a matter of hours. Switzerland, southern Germany, and Austria have a warmer climate due to the Foehn, as moist winds off the Mediterranean Sea blow over the Alps.
Etymology
The name Foehn (, ) arose in the Alpine region. Originating from Latin , a mild west wind of which Favonius was the Roman personification and probably transmitted by or just , the term was adopted as . In the Southern Alps, the phenomenon is known as but also and in Serbo-Croatian and Slovene. The German word (pronounced the same way) also means 'hairdryer', while the word is a genericized trademark today owned by AEG. The form phon is used in French-speaking parts of Switzerland as well as in Italy.
The name was originally used to refer to the south wind which blows during the winter months and brings thaw conditions to the northern side of the Alps. Because Föhn later became a generic term that was extended to other mountain ranges around the world that experience similar phenomena, the name "Alpine föhn" () was coined for the Föhns of the Alpine region.
Causes
There are four known causes of the Foehn warming and drying effect. These mechanisms often act together, with their contributions varying depending on the size and shape of the mountain barrier and on the meteorological conditions, such as the upstream wind speed, temperature and humidity.
Condensation and precipitation
When winds blow over elevated terrain, air forced upwards expands and cools due to the decrease in pressure with height. Since colder air can hold less water vapor, moisture condenses to form clouds and precipitates as rain or snow on the mountain's upwind slopes. The change of state from vapor to liquid water releases latent heat energy which heats the air, partially countering the cooling that occurs as the air rises. The subsequent removal of moisture as precipitation renders this heat gain by the air irreversible, leading to the warm, dry, Foehn conditions as the air descends in the mountain's lee. This mechanism has become a popular textbook example of atmospheric thermodynamics. However, the common occurrence of 'dry' Foehn events, where there is no precipitation, implies there must be other mechanisms.
Isentropic draw-down
Isentropic draw-down is the draw-down of warmer, drier air from aloft. When the approaching winds are insufficiently strong to propel the low-level air up and over the mountain barrier, the airflow is said to be 'blocked' by the mountain and only air higher up near mountain-top level is able to pass over and down the lee slopes as Foehn winds. These higher source regions provide Foehn air that becomes warmer and drier on the leeside after it is compressed with descent due to the increase in pressure towards the surface.
Mechanical mixing
When river water passes over rocks, turbulence is generated in the form of rapids, and white water reveals the turbulent mixing of the water with the air above. Similarly, as air passes over mountains, turbulence occurs and the atmosphere is mixed in the vertical. This mixing generally leads to a downward warming and upward moistening of the cross-mountain airflow, and consequently to warmer, drier Foehn winds in the valleys downwind.
Radiative warming
Dry Foehn conditions are responsible for the occurrence of rain shadows in the lee of mountains, where clear, sunny conditions prevail. This often leads to greater daytime radiative (solar) warming under Foehn conditions. This type of warming is particularly important in cold regions where snow or ice melt is a concern or where avalanches are a risk.
Effects
Winds of this type are also called "snow-eaters" for their ability to make snow and ice melt or sublimate rapidly. This is a result not only of the warmth of Foehn air, but also its low relative humidity. Accordingly, Foehn winds are known to contribute to the disintegration of ice shelves in the polar regions.
Foehn winds are notorious among mountaineers in the Alps, especially those climbing the Eiger, for whom the winds add further difficulty in ascending an already difficult peak.
They are also associated with the rapid spread of wildfires, making some regions which experience these winds particularly fire-prone.
Purported physiological effects
Anecdotally, residents in areas of frequent Foehn winds have reported experiencing a variety of illnesses ranging from migraines to psychosis. The first clinical review of these effects was published by the Austrian physician Anton Czermak in the 19th century. A study by the Ludwig-Maximilians-Universität München found that suicide and accidents increased by 10 percent during Foehn winds in Central Europe. The causation of Föhnkrankheit (English: Foehn-sickness) is unproven. Labels for preparations of aspirin combined with caffeine, codeine and the like will sometimes include Föhnkrankheit among the indications. Evidence for effects from Chinook winds remains anecdotal, as it does for New Zealand's Nor'wester.
In some regions, Foehn winds are associated with causing circulatory problems, headaches, or similar ailments. Researchers have found, however, the Foehn wind's warm temperature to be beneficial to humans in most situations, and have theorized that the reported negative effects may be a result of secondary factors, such as changes in the electrical field or in the ion state of the atmosphere, the wind's relatively low humidity, or the generally unpleasant sensation of being in an environment with strong and gusty winds.
Local examples
Regionally, these winds are known by many different names. These include:
in Africa
Bergwind in South Africa
in the Americas
The Brookings Effect on the southwestern coast of Oregon, also known as the Chetco Effect.
Chinook winds east of the Rocky Mountains and the Cascade Range in the United States and Canada, and north, east and west of the Chugach Mountains of Alaska, United States
Foehn winds in the foothills of the southern Appalachian Mountains, which can be unusual compared to other Foehn winds in that the relative humidity typically changes little due to the increased moisture in the source air mass
The Santa Ana winds of southern California, including the Sundowner winds of Santa Barbara, are in some ways similar to the Föhn, but originate in dry deserts as a katabatic wind. However, traditional Föhn conditions frequently prevail along the Santa Monica and Santa Ana Mountains and their respective leeward valleys, the San Fernando Valley and the Riverside County portion of the Inland Empire region.
Puelche wind in Chile
Suêtes on the west coast of Cape Breton Island, Nova Scotia, Canada
Wreckhouse winds in the southwest corner of the island of Newfoundland, Newfoundland and Labrador, Canada
Zonda winds in Argentina
in Antarctica
Föhn wall on Signy Island, South Orkneys
in Asia
Garmesh, Garmij, Garmbaad (): (, ) in Gilan region (near the Alborz) in the south west of Caspian Sea in Iran.
In winter, a Foehn effect occurs in the West Azerbaijan province, Iran (around Lake Urmia) as manifested by the province's dry winters relative to those in the windward part of the region (Northern Iraq or Kurdistan Region and Hakkâri Province in Turkey). For example, the winter rainfall of Urmia and Salmas in Iranian Azerbaijan is much lower than Batifa and Soran in Iraqi Kurdistan, and Hakkâri in the Hakkâri Province, which are roughly on the same latitude but are on the windward side of the Zagros Mountains.
Loo in Indo-Gangetic Plain
Warm Braw in the Schouten Islands north of West Papua, Indonesia.
Wuhan in China is famously known as one of the Three Furnaces on account of its extremely hot weather in summer resulting from the adiabatic warming effect created by mountains further south.
Laos wind (), hot-dry west wind () in northern and central Vietnam.
in Europe
Favonio in Ticino and north-western Italy due to western and northern winds crossing the Alps (mostly in winter)
Garbino in the Adriatic coast of Italy due to south-western winds crossing the Apennine Mountains (mostly in fall and winter)
Fen in northwest Slovenia
Fønvind in South Norway, in particular Central Norway, resulting in extreme winter warming, including Scandinavia's warmest winter temperature in Sunndalsøra.
Fogony in the Catalan Pyrenees
Föhn or Foehn in Austria, southern Germany, Switzerland, France and Liechtenstein
Föhn in Ostrobothnia and Western Lapland in Finland as moist air crosses Scandinavian Mountains and dries up.
Halny in the Carpathian Mountains, southern Poland and northern Slovakia
The Helm Wind, on the Pennines in the Eden Valley, Cumbria, England
Hnjúkaþeyr in Icelandic
Lodos wind, causing warm temperatures in the leeward side of mountains in the mild-winter climate of the Aegean Sea, Greece and western Turkey, as well as unusually mild temperatures in the cool or moderately cold winter climates north of the Marmara Sea, such as Istanbul, Adapazarı and Zonguldak.
Košava (Koshava) wind in Serbia that blows along the Danube River
Nortada in Cascais, and most notoriously in Guincho Beach, making it one of the best windsurfing spots in Europe
Ponentà in Valencia (eastern Spain)
Terral in Málaga (southern Spain)
Viento del Sur (Southern Wind) or Hego haizea in Basque in the Cantabrian region (northern Spain)
in Oceania
The Great Dividing foehn in southeast Australia, leeward of the Great Dividing Range, observed in the coastal plains of New South Wales, and also in eastern Victoria and eastern Tasmania.
The Nor'wester in Hawkes Bay, Canterbury, and Otago, New Zealand
Gallery
| Physical sciences | Winds | Earth science |
351583 | https://en.wikipedia.org/wiki/Liquid%20helium | Liquid helium | Liquid helium is a physical state of helium at very low temperatures at standard atmospheric pressures. Liquid helium may show superfluidity.
At standard pressure, the chemical element helium exists in a liquid form only at the extremely low temperature of . Its boiling point and critical point depend on the isotope of helium present: the common isotope helium-4 or the rare isotope helium-3. These are the only two stable isotopes of helium. See the table below for the values of these physical quantities. The density of liquid helium-4 at its boiling point and a pressure of one atmosphere (101.3 kilopascals) is about , or about one-eighth the density of liquid water.
Liquefaction
Helium was first liquefied on July 10, 1908, by the Dutch physicist Heike Kamerlingh Onnes at the University of Leiden in the Netherlands. At that time, helium-3 was unknown because the mass spectrometer had not yet been invented. In more recent decades, liquid helium has been used as a cryogenic refrigerant (which is used in cryocoolers), and liquid helium is produced commercially for use in superconducting magnets such as those used in magnetic resonance imaging (MRI), nuclear magnetic resonance (NMR), magnetoencephalography (MEG), and experiments in physics, such as low temperature Mössbauer spectroscopy. The Large Hadron Collider contains superconducting magnets that are cooled with 120 tonnes of liquid helium.
Liquified helium-3
A helium-3 atom is a fermion and at very low temperatures, they form two-atom Cooper pairs which are bosonic and condense into a superfluid. These Cooper pairs are substantially larger than the interatomic separation.
Characteristics
The temperature required to produce liquid helium is low because of the weakness of the attractions between the helium atoms. These interatomic forces in helium are weak to begin with because helium is a noble gas, but the interatomic attractions are reduced even more by the effects of quantum mechanics. These are significant in helium because of its low atomic mass of about four atomic mass units. The zero point energy of liquid helium is less if its atoms are less confined by their neighbors. Hence in liquid helium, its ground state energy can decrease by a naturally occurring increase in its average interatomic distance. However at greater distances, the effects of the interatomic forces in helium are even weaker.
Because of the very weak interatomic forces in helium, the element remains a liquid at atmospheric pressure all the way from its liquefaction point down to absolute zero. At temperatures below their liquefaction points, both helium-4 and helium-3 undergo transitions to superfluids. (See the table below.) Liquid helium can be solidified only under very low temperatures and high pressures.
Liquid helium-4 and the rare helium-3 are not completely miscible. Below 0.9 kelvin at their saturated vapor pressure, a mixture of the two isotopes undergoes a phase separation into a normal fluid (mostly helium-3) that floats on a denser superfluid consisting mostly of helium-4. This phase separation happens because the overall mass of liquid helium can reduce its thermodynamic enthalpy by separating.
At extremely low temperatures, the superfluid phase, rich in helium-4, can contain up to 6% helium-3 in solution. This makes the small-scale use of the dilution refrigerator possible, which is capable of reaching temperatures of a few millikelvins.
Superfluid helium-4 has substantially different properties from ordinary liquid helium.
History
In 1908, Kamerlingh-Onnes succeeded in liquifying a small quantity of helium. In 1923, he provided advice to the Canadian physicist John Cunningham McLennan, who was the first to produce quantities of liquid helium almost on demand.
In 1932 Einstein reported that the liquid helium could help in creating an atomic bomb.
Important early work on the characteristics of liquid helium was done by the Soviet physicist Lev Landau, later extended by the American physicist Richard Feynman.
In 1961, Vignos and Fairbank reported the existence of a different phase of solid helium-4, designated the gamma-phase. It exists for a narrow range of pressure between 1.45 and 1.78 K.
Data
Gallery
| Physical sciences | s-Block | Chemistry |
351603 | https://en.wikipedia.org/wiki/Piloting | Piloting | Piloting or pilotage is the process of navigating on water or in the air using fixed points of reference on the sea or on land, usually with reference to a nautical chart or aeronautical chart to obtain a fix of the position of the vessel or aircraft with respect to a desired course or location. Horizontal fixes of position from known reference points may be obtained by sight or by radar. Vertical position may be obtained by depth sounder to determine depth of the water body below a vessel or by altimeter to determine an aircraft's altitude, from which its distance above the ground can be deduced. Piloting a vessel is usually practiced close to shore or on inland waterways. Pilotage of an aircraft is practiced under visual meteorological conditions for flight.
Land navigation is a related discipline, using a topographic map, especially when applied over trackless terrain. Divers use related techniques for underwater navigation.
Piloting references
Charts
Depending on whether one is navigating on a water course, in the air or on land, a different chart applies for the navigator:
Nautical charts – show coastal regions and depict depths of water and land features, natural features of the seabed, details of the coastline, navigational hazards, locations of natural and human-made aids to navigation, and human-made structures such as harbours, buildings and bridges.
Aeronautical charts – for visual meteorological conditions depict terrain, geographic features, navigational aids and other aids to navigation. They vary in scale from 1:1,000,000 for world aeronautical charts to 1:250,000.
Topographic maps – show landforms and terrain, lakes and rivers, forest cover, administrative areas, populated areas, roads and railways, and other man-made features.
Maritime piloting
Coastal mariners often use reference manuals, called "pilots" for navigating coastal waters. In addition to providing descriptions of shipping channels and coastal profiles, they discuss weather, currents and other topics of interest to mariners. Notable guides include a worldwide series of "Sailing Directions" by the United Kingdom Hydrographic Office (formerly by the British Admiralty) that includes, most notably, the English Channel, the Mediterranean Sea, the Red Sea and the Persian Gulf. Another series worldwide series of Sailing Directions is by the US National Geospatial-Intelligence Agency, which has planning guide and enroute portions. The "United States Coast Pilot", by the National Oceanic and Atmospheric Administration (NOAA) Office of Coast Survey, covers the coastal and intracoastal waters and the Great Lakes of the United States.
Points of reference
Common types of visual reference point used for piloting and pilotage include:
Day
Natural features: Mountains, hills, lakes, rivers and coastal features such as cliffs, rocks and beaches
Navigational aids: sea marks (including buoys and beacons) and landmarks
Other structures: Airports, cities, dams, highways, and radio antennas
Night
Lighted navigational aids: Lighthouses, lightvessels and lighted sea marks
Lighted structures: Airports, illuminated towers and buildings
Vertical
Depth, measured with a depth sounder or lead line, can be used to identify a bathymetric contour or crossing point. Similarly, elevation can be used to confirm a geographic contour or crossing point. Measurement of depth and altitude allow vessels and aircraft navigators to confirm clear passage over obstructions.
Fix of position
Instruments used
On shipboard, navigators may use a pelorus to obtain bearings, relative to the vessel, from charted objects. A hand bearing compass provides magnetic bearings. On land a hand compass provides bearings to landmarks.
Afloat
Mariners use position-fixing navigation, to obtain a "position fix" or "fix" by measuring the bearing of the navigator's current position from known points of reference. A visual fix of position can be made by using any sighting device with a bearing indicator to obtain position lines from the navigator's current position to each point of reference. Two or more objects of known position are sighted as points of reference, and the bearings recorded. Bearing lines or transits are then plotted on a chart through the locations of the sighted items. The intersection of these lines is then the current position of the navigator.
Usually, a fix is where two or more position lines intersect at any given time. If three position lines can be obtained, the resulting "cocked hat", where the 3 lines do not intersect at the same point, but create a triangle where the vessel is inside, gives the navigator an indication of the accuracy in the three separate position lines.
If two geographic features are visually aligned (the edge of an island aligned with the edge of an island behind, a flag pole and a building, etc.), the extension of the line joining the features is called a "transit". A transit is not affected by compass accuracy, and is often used to check a compass for errors.
The most accurate fixes occur when the position lines are at right angles to each other.
Aloft
Flying at low altitudes and with sufficient visibility, aircraft pilots use nearby rivers, roads, railroad tracks and other visual references to establish their position.
Course versus ground track
The line connecting fixes is the track over the ground or sea bottom. The navigator compares the ground track with the navigational course for that leg of the intended route, in order to make a correction in "heading", the direction in which the craft is pointed to maintain its course in compensation for cross-currents of wind or water that may carry the craft off course.
In channels and rivers
Where a channel is narrow, as in some harbor entrances and on some rivers, a system of beacons allows mariners to align pairs of daymarks, called "range markers", to form a "leading line" (British English) or "range axis" (American English), along which to navigate safely. When lighted, these markers are called "leading lights" (British English) or "range lights" (American English). The relative positions of the marks and the vessel affect the accuracy of perceiving the leading line.
| Technology | Navigation | null |
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