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1,602,252 | https://en.wikipedia.org/wiki/Source%20counts | The source counts distribution of radio-sources from a radio-astronomical survey is the cumulative distribution of the number of sources (N) brighter than a given flux density (S). As it is usually plotted on a log-log scale its distribution is known as the log N – log S plot. It is one of several cosmological tests that were conceived in the 1930s to check the viability of and compare new cosmological models.
Early work to catalogue radio sources aimed to determine the source count distribution as a discriminating test of different cosmological models. For example, a uniform distribution of radio sources at low redshift, such as might be found in a 'steady-state Euclidean universe,' would produce a slope of −1.5 in the cumulative distribution of log(N) versus log(S).
Data from the early Cambridge 2C survey (published 1955) apparently implied a (log(N), log(S)) slope of nearly −3.0. This appeared to invalidate the steady state theory of Fred Hoyle, Hermann Bondi and Thomas Gold. Unfortunately many of these weaker sources were subsequently found to be due to 'confusion' (the blending of several weak sources in the side-lobes of the interferometer, producing a stronger response).
By contrast, analysis from the contemporaneous Mills Cross data (by Slee and Mills) were consistent with an index of −1.5.
Later and more accurate surveys from Cambridge, 3C, 3CR, and 4C, also showed source count slopes steeper than −1.5, though by a smaller margin than 2C. This convinced some cosmologists that the steady state theory was wrong, although residual problems with confusion provided some defense for Hoyle and his colleagues.
The immediate interest in testing the steady-state theory through source-counts was reduced by the discovery of the 3K microwave background radiation in the mid-1960s, which essentially confirmed the Big-Bang model.
Later radio survey data have shown a complex picture — the 3C and 4C claims appear to hold up, while at fainter levels the source counts flatten substantially below a slope of −1.5. This is now understood to reflect the effects of both density and luminosity evolution of the principal radio sources over cosmic timescales.
See also
Tolman surface brightness test
References
Physical cosmology | Source counts | [
"Physics",
"Astronomy"
] | 488 | [
"Astronomical sub-disciplines",
"Theoretical physics",
"Physical cosmology",
"Astrophysics"
] |
1,602,490 | https://en.wikipedia.org/wiki/Extensive-form%20game | In game theory, an extensive-form game is a specification of a game allowing (as the name suggests) for the explicit representation of a number of key aspects, like the sequencing of players' possible moves, their choices at every decision point, the (possibly imperfect) information each player has about the other player's moves when they make a decision, and their payoffs for all possible game outcomes. Extensive-form games also allow for the representation of incomplete information in the form of chance events modeled as "moves by nature". Extensive-form representations differ from normal-form in that they provide a more complete description of the game in question, whereas normal-form simply boils down the game into a payoff matrix.
Finite extensive-form games
Some authors, particularly in introductory textbooks, initially define the extensive-form game as being just a game tree with payoffs (no imperfect or incomplete information), and add the other elements in subsequent chapters as refinements. Whereas the rest of this article follows this gentle approach with motivating examples, we present upfront the finite extensive-form games as (ultimately) constructed here. This general definition was introduced by Harold W. Kuhn in 1953, who extended an earlier definition of von Neumann from 1928. Following the presentation from , an n-player extensive-form game thus consists of the following:
A finite set of n (rational) players
A rooted tree, called the game tree
Each terminal (leaf) node of the game tree has an n-tuple of payoffs, meaning there is one payoff for each player at the end of every possible play
A partition of the non-terminal nodes of the game tree in n+1 subsets, one for each (rational) player, and with a special subset for a fictitious player called Chance (or Nature). Each player's subset of nodes is referred to as the "nodes of the player". (A game of complete information thus has an empty set of Chance nodes.)
Each node of the Chance player has a probability distribution over its outgoing edges.
Each set of nodes of a rational player is further partitioned in information sets, which make certain choices indistinguishable for the player when making a move, in the sense that:
there is a one-to-one correspondence between outgoing edges of any two nodes of the same information set—thus the set of all outgoing edges of an information set is partitioned in equivalence classes, each class representing a possible choice for a player's move at some point—, and
every (directed) path in the tree from the root to a terminal node can cross each information set at most once
the complete description of the game specified by the above parameters is common knowledge among the players
A play is thus a path through the tree from the root to a terminal node. At any given non-terminal node belonging to Chance, an outgoing branch is chosen according to the probability distribution. At any rational player's node, the player must choose one of the equivalence classes for the edges, which determines precisely one outgoing edge except (in general) the player doesn't know which one is being followed. (An outside observer knowing every other player's choices up to that point, and the realization of Nature's moves, can determine the edge precisely.) A pure strategy for a player thus consists of a selection—choosing precisely one class of outgoing edges for every information set (of his). In a game of perfect information, the information sets are singletons. It's less evident how payoffs should be interpreted in games with Chance nodes. It is assumed that each player has a von Neumann–Morgenstern utility function defined for every game outcome; this assumption entails that every rational player will evaluate an a priori random outcome by its expected utility.
The above presentation, while precisely defining the mathematical structure over which the game is played, elides however the more technical discussion of formalizing statements about how the game is played like "a player cannot distinguish between nodes in the same information set when making a decision". These can be made precise using epistemic modal logic; see for details.
A perfect information two-player game over a game tree (as defined in combinatorial game theory and artificial intelligence) can be represented as an extensive form game with outcomes (i.e. win, lose, or draw). Examples of such games include tic-tac-toe, chess, and infinite chess. A game over an expectminimax tree, like that of backgammon, has no imperfect information (all information sets are singletons) but has moves of chance. For example, poker has both moves of chance (the cards being dealt) and imperfect information (the cards secretly held by other players).
Perfect and complete information
A complete extensive-form representation specifies:
the players of a game
for every player every opportunity they have to move
what each player can do at each of their moves
what each player knows for every move
the payoffs received by every player for every possible combination of moves
The game on the right has two players: 1 and 2. The numbers by every non-terminal node indicate to which player that decision node belongs. The numbers by every terminal node represent the payoffs to the players (e.g. 2,1 represents a payoff of 2 to player 1 and a payoff of 1 to player 2). The labels by every edge of the graph are the name of the action that edge represents.
The initial node belongs to player 1, indicating that player 1 moves first. Play according to the tree is as follows: player 1 chooses between U and D; player 2 observes player 1's choice and then chooses between U' and D' . The payoffs are as specified in the tree. There are four outcomes represented by the four terminal nodes of the tree: (U,U'), (U,D'), (D,U') and (D,D'). The payoffs associated with each outcome respectively are as follows (0,0), (2,1), (1,2) and (3,1).
If player 1 plays D, player 2 will play U' to maximise their payoff and so player 1 will only receive 1. However, if player 1 plays U, player 2 maximises their payoff by playing D' and player 1 receives 2. Player 1 prefers 2 to 1 and so will play U and player 2 will play D' . This is the subgame perfect equilibrium.
Imperfect information
An advantage of representing the game in this way is that it is clear what the order of play is. The tree shows clearly that player 1 moves first and player 2 observes this move. However, in some games play does not occur like this. One player does not always observe the choice of another (for example, moves may be simultaneous or a move may be hidden). An information set is a set of decision nodes such that:
Every node in the set belongs to one player.
When the game reaches the information set, the player who is about to move cannot differentiate between nodes within the information set; i.e. if the information set contains more than one node, the player to whom that set belongs does not know which node in the set has been reached.
In extensive form, an information set is indicated by a dotted line connecting all nodes in that set or sometimes by a loop drawn around all the nodes in that set.
If a game has an information set with more than one member that game is said to have imperfect information. A game with perfect information is such that at any stage of the game, every player knows exactly what has taken place earlier in the game; i.e. every information set is a singleton set. Any game without perfect information has imperfect information.
The game on the right is the same as the above game except that player 2 does not know what player 1 does when they come to play. The first game described has perfect information; the game on the right does not. If both players are rational and both know that both players are rational and everything that is known by any player is known to be known by every player (i.e. player 1 knows player 2 knows that player 1 is rational and player 2 knows this, etc. ad infinitum), play in the first game will be as follows: player 1 knows that if they play U, player 2 will play D' (because for player 2 a payoff of 1 is preferable to a payoff of 0) and so player 1 will receive 2. However, if player 1 plays D, player 2 will play U' (because to player 2 a payoff of 2 is better than a payoff of 1) and player 1 will receive 1. Hence, in the first game, the equilibrium will be (U, D' ) because player 1 prefers to receive 2 to 1 and so will play U and so player 2 will play D' .
In the second game it is less clear: player 2 cannot observe player 1's move. Player 1 would like to fool player 2 into thinking they have played U when they have actually played D so that player 2 will play D' and player 1 will receive 3. In fact in the second game there is a perfect Bayesian equilibrium where player 1 plays D and player 2 plays U' and player 2 holds the belief that player 1 will definitely play D. In this equilibrium, every strategy is rational given the beliefs held and every belief is consistent with the strategies played. Notice how the imperfection of information changes the outcome of the game.
To more easily solve this game for the Nash equilibrium, it can be converted to the normal form. Given this is a simultaneous/sequential game, player one and player two each have two strategies.
Player 1's Strategies: {U , D}
Player 2's Strategies: {U’ , D’}
We will have a two by two matrix with a unique payoff for each combination of moves. Using the normal form game, it is now possible to solve the game and identify dominant strategies for both players.
If player 1 plays Up (U), player 2 prefers to play Down (D’) (Payoff 1>0)
If player 1 plays Down (D), player 2 prefers to play Up (U’) (Payoff 2>1)
If player 2 plays Up (U’), player 1 prefers to play Down (D) (Payoff 1>0)
If player 2 plays Down (D’), player 1 prefers to play Down (D) (3>2)
These preferences can be marked within the matrix, and any box where both players have a preference provides a nash equilibrium. This particular game has a single solution of (D,U’) with a payoff of (1,2).
In games with infinite action spaces and imperfect information, non-singleton information sets are represented, if necessary, by inserting a dotted line connecting the (non-nodal) endpoints behind the arc described above or by dashing the arc itself. In the Stackelberg competition described above, if the second player had not observed the first player's move the game would no longer fit the Stackelberg model; it would be Cournot competition.
Incomplete information
It may be the case that a player does not know exactly what the payoffs of the game are or of what type their opponents are. This sort of game has incomplete information. In extensive form it is represented as a game with complete but imperfect information using the so-called Harsanyi transformation. This transformation introduces to the game the notion of nature's choice or God's choice. Consider a game consisting of an employer considering whether to hire a job applicant. The job applicant's ability might be one of two things: high or low. Their ability level is random; they either have low ability with probability 1/3 or high ability with probability 2/3. In this case, it is convenient to model nature as another player of sorts who chooses the applicant's ability according to those probabilities. Nature however does not have any payoffs. Nature's choice is represented in the game tree by a non-filled node. Edges coming from a nature's choice node are labelled with the probability of the event it represents occurring.
The game on the left is one of complete information (all the players and payoffs are known to everyone) but of imperfect information (the employer doesn't know what nature's move was.) The initial node is in the centre and it is not filled, so nature moves first. Nature selects with the same probability the type of player 1 (which in this game is tantamount to selecting the payoffs in the subgame played), either t1 or t2. Player 1 has distinct information sets for these; i.e. player 1 knows what type they are (this need not be the case). However, player 2 does not observe nature's choice. They do not know the type of player 1; however, in this game they do observe player 1's actions; i.e. there is perfect information. Indeed, it is now appropriate to alter the above definition of complete information: at every stage in the game, every player knows what has been played by the other players. In the case of private information, every player knows what has been played by nature. Information sets are represented as before by broken lines.
In this game, if nature selects t1 as player 1's type, the game played will be like the very first game described, except that player 2 does not know it (and the very fact that this cuts through their information sets disqualify it from subgame status). There is one separating perfect Bayesian equilibrium; i.e. an equilibrium in which different types do different things.
If both types play the same action (pooling), an equilibrium cannot be sustained. If both play D, player 2 can only form the belief that they are on either node in the information set with probability 1/2 (because this is the chance of seeing either type). Player 2 maximises their payoff by playing D' . However, if they play D' , type 2 would prefer to play U. This cannot be an equilibrium. If both types play U, player 2 again forms the belief that they are at either node with probability 1/2. In this case player 2 plays D' , but then type 1 prefers to play D.
If type 1 plays U and type 2 plays D, player 2 will play D' whatever action they observe, but then type 1 prefers D. The only equilibrium hence is with type 1 playing D, type 2 playing U and player 2 playing U' if they observe D and randomising if they observe U. Through their actions, player 1 has signalled their type to player 2.
Formal definition
Formally, a finite game in extensive form is a structure
where:
is a finite tree with a set of nodes , a unique initial node , a set of terminal nodes (let be a set of decision nodes) and an immediate predecessor function on which the rules of the game are represented,
is a partition of called an information partition,
is a set of actions available for each information set which forms a partition on the set of all actions .
is an action partition associating each node to a single action , fulfilling:
, the restriction of on is a bijection, with the set of successor nodes of .
is a finite set of players, is (a special player called) nature, and is a player partition of information set . Let be a single player that makes a move at node .
is a family of probabilities of the actions of nature, and
is a payoff profile function.
Infinite action space
It may be that a player has an infinite number of possible actions to choose from at a particular decision node. The device used to represent this is an arc joining two edges protruding from the decision node in question. If the action space is a continuum between two numbers, the lower and upper delimiting numbers are placed at the bottom and top of the arc respectively, usually with a variable that is used to express the payoffs. The infinite number of decision nodes that could result are represented by a single node placed in the centre of the arc. A similar device is used to represent action spaces that, whilst not infinite, are large enough to prove impractical to represent with an edge for each action.
The tree on the left represents such a game, either with infinite action spaces (any real number between 0 and 5000) or with very large action spaces (perhaps any integer between 0 and 5000). This would be specified elsewhere. Here, it will be supposed that it is the former and, for concreteness, it will be supposed it represents two firms engaged in Stackelberg competition. The payoffs to the firms are represented on the left, with and as the strategy they adopt and and as some constants (here marginal costs to each firm). The subgame perfect Nash equilibria of this game can be found by taking the first partial derivative of each payoff function with respect to the follower's (firm 2) strategy variable () and finding its best response function, . The same process can be done for the leader except that in calculating its profit, it knows that firm 2 will play the above response and so this can be substituted into its maximisation problem. It can then solve for by taking the first derivative, yielding . Feeding this into firm 2's best response function, and is the subgame perfect Nash equilibrium.
See also
Axiom of determinacy
Perfect information
Combinatorial game theory
Self-confirming equilibrium
Sequential game
Signalling
Solution concept
References
Dresher M. (1961). The mathematics of games of strategy: theory and applications (Ch4: Games in extensive form, pp74–78). Rand Corp.
Fudenberg D and Tirole J. (1991) Game theory (Ch3 Extensive form games, pp67–106). MIT press.
. An 88-page mathematical introduction; see Chapters 4 and 5. Free online at many universities.
Luce R.D. and Raiffa H. (1957). Games and decisions: introduction and critical survey. (Ch3: Extensive and Normal Forms, pp39–55). Wiley New York.
Osborne MJ and Rubinstein A. 1994. A course in game theory (Ch6 Extensive game with perfect information, pp. 89–115). MIT press.
. A comprehensive reference from a computational perspective; see Chapter 5. Downloadable free online.
Further reading
, 6.1, "Disasters in Game Theory" and 7.2 "Measurability (The Axiom of Determinateness)", discusses problems in extending the finite-case definition to infinite number of options (or moves)
Historical papers
contains Kuhn's lectures at Princeton from 1952 (officially unpublished previously, but in circulation as photocopies)
Game theory game classes | Extensive-form game | [
"Mathematics"
] | 3,891 | [
"Game theory game classes",
"Game theory"
] |
1,603,001 | https://en.wikipedia.org/wiki/Complete%20information | In economics and game theory, complete information is an economic situation or game in which knowledge about other market participants or players is available to all participants. The utility functions (including risk aversion), payoffs, strategies and "types" of players are thus common knowledge. Complete information is the concept that each player in the game is aware of the sequence, strategies, and payoffs throughout gameplay. Given this information, the players have the ability to plan accordingly based on the information to maximize their own strategies and utility at the end of the game. A typical example is the prisoner's dilemma.
Inversely, in a game with incomplete information, players do not possess full information about their opponents. Some players possess private information, a fact that the others should take into account when forming expectations about how those players will behave. A typical example is an auction: each player knows their own utility function (valuation for the item), but does not know the utility function of the other players.
Applications
Games of incomplete information arise frequently in social science. For instance, John Harsanyi was motivated by consideration of arms control negotiations, where the players may be uncertain both of the capabilities of their opponents and of their desires and beliefs.
It is often assumed that the players have some statistical information about the other players, e.g. in an auction, each player knows that the valuations of the other players are drawn from some probability distribution. In this case, the game is called a Bayesian game.
In games that have a varying degree of complete information and game type, there are different methods available to the player to solve the game based on this information. In games with static, complete information, the approach to solve is to use Nash equilibrium to find viable strategies. In dynamic games with complete information, backward induction is the solution concept, which eliminates non-credible threats as potential strategies for players.
A classic example of a dynamic game with complete information is Stackelberg's (1934) sequential-move version of Cournot duopoly. Other examples include Leontief's (1946) monopoly-union model and Rubenstein's bargaining model.
Lastly, when complete information is unavailable (incomplete information games), these solutions turn towards Bayesian Nash Equilibria since games with incomplete information become Bayesian games. In a game of complete information, the players' payoffs functions are common knowledge, whereas in a game of incomplete information at least one player is uncertain about another player's payoff function.
Extensive form
The extensive form can be used to visualize the concept of complete information. By definition, players know where they are as depicted by the nodes, and the final outcomes as illustrated by the utility payoffs. The players also understand the potential strategies of each player and as a result their own best course of action to maximize their payoffs.
Complete versus perfect information
Complete information is importantly different from perfect information.
In a game of complete information, the structure of the game and the payoff functions of the players are commonly known but players may not see all of the moves made by other players (for instance, the initial placement of ships in Battleship); there may also be a chance element (as in most card games). Conversely, in games of perfect information, every player observes other players' moves, but may lack some information on others' payoffs, or on the structure of the game. A game with complete information may or may not have perfect information, and vice versa.
Examples of games with imperfect but complete information are card games, where each player's cards are hidden from other players but objectives are known, as in contract bridge and poker, if the outcomes are assumed to be binary (players can only win or lose in a zero-sum game). Games with complete information generally require one player to outwit the other by forcing them to make risky assumptions.
Examples of games with incomplete but perfect information are conceptually more difficult to imagine, such as a Bayesian game. A game of chess is a commonly given example to illustrate how the lack of certain information influences the game, without chess itself being such a game. One can readily observe all of the opponent's moves and viable strategies available to them but never ascertain which one the opponent is following until this might prove disastrous for one. Games with perfect information generally require one player to outwit the other by making them misinterpret one's decisions.
See also
Bayesian game
Handicap principle
Market impact
Screening game
Signaling game
Small talk
Trash-talk
References
Bibliography
Watson, J. (2015) Strategy: An Introduction to Game Theory. Volume 139. New York, WW Norton
Fudenberg, D. and Tirole, J. (1993) Game Theory. MIT Press. (see Chapter 6, sect 1)
Gibbons, R. (1992) A primer in game theory. Harvester-Wheatsheaf. (see Chapter 3)
Ian Frank, David Basin (1997), Artificial Intelligence 100 (1998) 87-123. "Search in games with incomplete information: a case study using Bridge card play".
Game theory
Perfect competition | Complete information | [
"Mathematics"
] | 1,044 | [
"Game theory"
] |
1,603,459 | https://en.wikipedia.org/wiki/Mittag-Leffler%20function | In mathematics, the Mittag-Leffler functions are a family of special functions. They are complex-valued functions of a complex argument z, and moreover depend on one or two complex parameters.
The one-parameter Mittag-Leffler function, introduced by Gösta Mittag-Leffler in 1903,
can be defined by the Maclaurin series
where is the gamma function, and is a complex parameter with .
The two-parameter Mittag-Leffler function, introduced by Wiman in 1905, is occasionally called the generalized Mittag-Leffler function. It has an additional complex parameter , and may be defined by the series
When , the one-parameter function is recovered.
In the case and are real and positive, the series converges for all values of the argument , so the Mittag-Leffler function is an entire function. This class of functions are important in the theory of the fractional calculus.
See below for three-parameter generalizations.
Some basic properties
For , the Mittag-Leffler function is an entire function of order , and type for any value of . In some sense, the Mittag-Leffler function is the simplest entire function of its order. The indicator function of is
This result actually holds for as well with some restrictions on when .
The Mittag-Leffler function satisfies the recurrence property (Theorem 5.1 of )
from which the following asymptotic expansion holds : for and real such that
then for all , we can show the following asymptotic expansions (Section 6. of ):
-as :
,
-and as :
.
A simpler estimate that can often be useful is given, thanks to the fact that the order and type of is and , respectively:
for any positive and any .
Special cases
For , the series above equals the Taylor expansion of the geometric series and consequently .
For we find: (Section 2 of )
Error function:
Exponential function:
Hyperbolic cosine:
For , we have
For , the integral
gives, respectively: , , .
Mittag-Leffler's integral representation
The integral representation of the Mittag-Leffler function is (Section 6 of )
where the contour starts and ends at and circles around the singularities and branch points of the integrand.
Related to the Laplace transform and Mittag-Leffler summation is the expression (Eq (7.5) of with )
Three-parameter generalizations
One generalization, characterized by three parameters, is
where and are complex parameters and .
Another generalization is the Prabhakar function
where is the Pochhammer symbol.
Applications of Mittag-Leffler function
One of the applications of the Mittag-Leffler function is in modeling fractional order viscoelastic materials. Experimental investigations into the time-dependent relaxation behavior of viscoelastic materials are characterized by a very fast decrease of the stress at the beginning of the relaxation process and an extremely slow decay for large times. It can even take a long time before a constant asymptotic value is reached. Therefore, a lot of Maxwell elements are required to describe relaxation behavior with sufficient accuracy. This ends in a difficult optimization problem in order to identify a large number of material parameters. On the other hand, over the years, the concept of fractional derivatives has been introduced to the theory of viscoelasticity. Among these models, the fractional Zener model was found to be very effective to predict the dynamic nature of rubber-like materials with only a small number of material parameters. The solution of the corresponding constitutive equation leads to a relaxation function of the Mittag-Leffler type. It is defined by the power series with negative arguments. This function represents all essential properties of the relaxation process under the influence of an arbitrary and continuous signal with a jump at the origin.
See also
Mittag-Leffler summation
Mittag-Leffler distribution
Notes
R Package 'MittagLeffleR' by Gurtek Gill, Peter Straka. Implements the Mittag-Leffler function, distribution, random variate generation, and estimation.
References
Gorenflo R., Kilbas A.A., Mainardi F., Rogosin S.V., Mittag-Leffler Functions, Related Topics and Applications (Springer, New York, 2014) 443 pages
External links
Mittag-Leffler function: MATLAB code
Mittag-Leffler and stable random numbers: Continuous-time random walks and stochastic solution of space-time fractional diffusion equations
Special functions
Analytic functions | Mittag-Leffler function | [
"Mathematics"
] | 952 | [
"Special functions",
"Combinatorics"
] |
1,603,488 | https://en.wikipedia.org/wiki/Elevator%20pitch | An elevator pitch, elevator speech, lift speech, or elevator statement is a short description of an idea, product, or company that explains the concept in a way such that any listener can understand it in a short period of time. This description typically explains who the thing is for, what it does, why it is needed, and how it will get done. When explaining an individual person, the description generally explains one's skills and goals, and why they would be a productive and beneficial person to have on a team or within a company or project. An elevator pitch does not have to include all of these components, but it usually does at least explain what the idea, product, company, or person is and their value.
Unlike a sales pitch, an elevator pitch can be used in a variety of ways, and may not have a clear buyer-seller relationship. The goal is simply to convey the overall concept or topic being pitched in an exciting way.
The name—elevator pitch—reflects the idea that it should be possible to deliver the summary in the time span of an elevator ride, or approximately thirty seconds to two minutes.
Background information and history
There are many origin stories for the elevator pitch. One commonly-known origin story is that of Ilene Rosenzweig and Michael Caruso, two former journalists active in the 1990s. According to Rosenzweig, Caruso was a senior editor at Vanity Fair and was continuously attempting to pitch story ideas to the Editor-In-Chief at the time, but could never pin her down long enough to do so simply because she was always on the move. So, in order to pitch her ideas, Caruso would join her during short free periods of time she had, such as on an elevator ride. Thus, the concept of an elevator pitch was created, as says Rosenzweig.
However, there is another known potential origin story that dates back before the story of Rosenzweig and Caruso. Philip Crosby, author of The Art of Getting Your Own Sweet Way (1972) and Quality Is Still Free (1996) suggested individuals should have a pre-prepared speech that can deliver information regarding themselves or a quality that they can provide within a short period of time, namely the amount of time of an elevator ride for if an individual finds themselves on an elevator with a prominent figure. Essentially, an elevator pitch is meant to allow an individual to, with very limited time, pitch themselves or an idea to a person who is high up in a company.
Crosby, who worked as a quality test technician, and then later as the Director of Quality at International Telephone and Telegraph, recounted how an elevator pitch could be used to push for change within the company. He planned a speech regarding the change he wanted to see and waited at the elevator at ITT headquarters. Crosby stepped onto an elevator with the CEO of the company to deliver his speech. Once they reached the floor where the CEO was getting off, Crosby was asked to deliver a full presentation on the topic at a meeting for all of the general managers.
Aspects
An elevator pitch is meant to last the duration of an elevator ride, which can vary in length from approximately thirty seconds to two minutes. Therefore, the main focus of an elevator pitch should be making it short and direct. According to the Idaho Business Review, the first two sentences of any elevator pitch are the most important, and should hook or grab the attention of the listener. Information in an elevator pitch, due to the limited amount of time, should be condensed to express the most important ideas or concepts within the allotted time.
The Idaho Business Review also suggests individuals who use an elevator pitch deliver it using simple language, avoiding statistics or other language that may disrupt the focus of the listener. Bloomberg Businessweek suggests that an important lesson to think about when giving an elevator pitch is to "adjust the pitch to the person who is listening, and refine it as you and your business continue to grow and change." When delivering an elevator pitch, individuals are encouraged to remain flexible and adaptable, and to be able to deliver the pitch in a genuine and fluent fashion. By doing so, the intended audience of the pitch will likely be able to follow the information, and will not consider it as being scripted.
Advantages
Advantages to conducting an elevator pitch include convenience and simplicity. For instance, elevator pitches can be given on short notice and without much preparation, making the listener more comfortable. Furthermore, elevator pitches allow the individual who is giving the pitch the ability to simplify the content and deliver it in a less complicated manner by providing the information in a cut-down fashion that gets right to the point.
See also
Business plan
High concept
Lightning talk
Mission statement
SWOT analysis
Vision statement
References
Further reading
Business terms
Elevators
Rhetorical techniques
Selling techniques
Statements | Elevator pitch | [
"Engineering"
] | 961 | [
"Building engineering",
"Elevators"
] |
1,603,557 | https://en.wikipedia.org/wiki/Cadmium%20poisoning | Cadmium is a naturally occurring toxic metal with common exposure in industrial workplaces, plant soils, and from smoking. Due to its low permissible exposure in humans, overexposure may occur even in situations where only trace quantities of cadmium are found. Cadmium is used extensively in electroplating, although the nature of the operation does not generally lead to overexposure. Cadmium is also found in some industrial paints and may represent a hazard when sprayed. Operations involving removal of cadmium paints by scraping or blasting may pose a significant hazard. The primary use of cadmium is in the manufacturing of NiCd rechargeable batteries. The primary source for cadmium is as a byproduct of refining zinc metal. Exposures to cadmium are addressed in specific standards for the general industry, shipyard employment, the construction industry, and the agricultural industry.
Signs and symptoms
Acute
Acute exposure to cadmium fumes may cause flu-like symptoms including chills, fever, and muscle ache sometimes referred to as "the cadmium blues." Symptoms may resolve after a week if there is no respiratory damage. More severe exposures can cause tracheobronchitis, pneumonitis, and pulmonary edema. Symptoms of inflammation may start hours after the exposure and include cough, dryness and irritation of the nose and throat, headache, dizziness, weakness, fever, chills, and chest pain.
Chronic
Complications of cadmium poisoning include cough, anemia, and kidney failure (possibly leading to death). Cadmium exposure increases one's chances of developing cancer. Similar to zinc, long-term exposure to cadmium fumes can cause lifelong anosmia.
Bone and joints
One of the main effects of cadmium poisoning is weak and brittle bones. The bones become soft (osteomalacia), lose bone mineral density (osteoporosis), and become weaker. This results in joint and back pain, and increases the risk of fractures. Spinal and leg pain is common, and a waddling gait often develops due to bone deformities caused by the long-term cadmium exposure. The pain eventually becomes debilitating, with fractures becoming more common as the bone weakens. Permanent deformation in bones can occur. In extreme cases of cadmium poisoning, mere body weight causes a fracture.
Renal
The kidney damage inflicted by cadmium poisoning is irreversible. The kidneys can shrink up to 30 percent. The kidneys lose their function to remove acids from the blood in proximal renal tubular dysfunction. The proximal renal tubular dysfunction causes hypophosphatemia, leading to muscle weakness and sometimes coma. Hyperchloremia also occurs. Kidney dysfunction also causes gout, a form of arthritis due to the accumulation of uric acid crystals in the joints because of high acidity of the blood (hyperuricemia). Cadmium exposure is also associated with the development of kidney stones.
Sources of exposure
Smoking is a significant source of cadmium exposure. Even small amounts of cadmium from smoking are highly toxic to humans, as the lungs absorb cadmium more efficiently than the stomach. Cadmium is emitted to the electronic cigarette (EC) aerosol but, on currently available data, the lifetime cancer risk (LCR) calculated does not exceed the acceptable risk limit.
Environmental
Buildup of cadmium levels in the water, air, and soil has been occurring particularly in industrial areas. Environmental exposure to cadmium has been particularly problematic in Japan where many people have consumed rice that was grown in cadmium-contaminated irrigation water. This phenomenon is known as itai-itai disease.
People who live near hazardous waste sites or factories that release cadmium into the air have the potential for exposure to cadmium in air. However, numerous state and federal regulations in the United States control the amount of cadmium that can be released to the air from waste sites and incinerators so that properly regulated sites are not hazardous. The general population and people living near hazardous waste sites may be exposed to cadmium in contaminated food, dust, or water from unregulated or accidental releases. Numerous regulations and use of pollution controls are enforced to prevent such releases.
Some sources of phosphate in fertilizers contain cadmium in amounts of up to 100 mg/kg, which can lead to an increase in the concentration of cadmium in soil (for example in New Zealand).
Food
Food is another source of cadmium. Plants may contain small or moderate amounts in non-industrial areas, but high levels may be found in the liver and kidneys of adult animals. The daily intake of cadmium through food varies by geographic region. Intake is reported to be approximately 8 to 30μg in Europe and the United States versus 59 to 113 μg in various areas of Japan. A small study of premium dark chocolate samples found 48% had high levels of cadmium, the source commonly being the presence of cadmium in soil in which they were grown.
Occupational exposure
In the 1950s and 1960s industrial exposure to cadmium was high, but as the toxic effects of cadmium became apparent, industrial limits on cadmium exposure have been reduced in most industrialized nations and many policy makers agree on the need to reduce exposure further. While working with cadmium it is important to do so under a fume hood to protect against dangerous fumes. Brazing fillers which contain cadmium should be handled with care. Serious toxicity problems have resulted from long-term exposure to cadmium plating baths.
Workers can be exposed to cadmium in air from the smelting and refining of metals, or from the air in plants that make cadmium products such as batteries, coatings, or plastics. Workers can also be exposed when soldering or welding metal that contains cadmium. Approximately 512,000 workers in the United States are in environments each year where cadmium exposure may occur. Regulations that set permissible levels of exposure, however, are enforced to protect workers and to make sure that levels of cadmium in the air are considerably below levels thought to result in harmful effects.
Artists who work with cadmium pigments, which are commonly used in strong oranges, reds, and yellows, can easily accidentally ingest dangerous amounts, particularly if they use the pigments in dry form, as with chalk pastels, or in mixing their own paints.
Consumer products
Cadmium is used in nickel-cadmium batteries; these are some of the most popular and most common cadmium-based products.
In February 2010, cadmium was found in an entire line of Wal-Mart exclusive Miley Cyrus jewelry. The charms were tested at the behest of the Associated Press and were found to contain high levels of cadmium. Wal-Mart did not stop selling the jewelry until May 12 because "it would be too difficult to test products already on its shelves".
On June 4, 2010, cadmium was detected in the paint used on promotional drinking glasses for the movie Shrek Forever After, sold by McDonald's Restaurants, triggering a recall of 12 million glasses.
Toxicology
Cadmium is an extremely toxic industrial and environmental pollutant classified as a human carcinogen: Group 1, according to the International Agency for Research on Cancer; Group 2a, according to Environmental Protection Agency (EPA); and a 1B carcinogen as classified by European Chemical Agency.
Toxicodynamics
Cellular toxicology
Inside cells, cadmium ions act as a catalytic hydrogen peroxide generator. This sudden surge of cytosolic hydrogen peroxide causes increased lipid peroxidation and additionally depletes ascorbate and glutathione stores. Hydrogen peroxide can also convert thiol groups on proteins into nonfunctional sulfonic acids and is also capable of directly attacking nuclear DNA. This oxidative stress causes the afflicted cell to manufacture large amounts of inflammatory cytokines.
Toxicokinetics
Inhaling cadmium-laden dust quickly leads to respiratory tract and kidney problems which can be fatal (often from kidney failure). Ingestion of any significant amount of cadmium causes immediate poisoning and damage to the liver and the kidneys. Compounds containing cadmium are also carcinogenic.
Diagnosis
Biomarkers of excessive exposure
Increased concentrations of urinary beta-2 microglobulin can be an early indicator of kidney dysfunction in persons chronically exposed to low but excessive levels of environmental cadmium. The urinary beta-2 microglobulin test is an indirect method of measuring cadmium exposure. Under some circumstances, the Occupational Health and Safety Administration requires screening for kidney damage in workers with long-term exposure to high levels of cadmium. Blood or urine cadmium concentrations provide a better index of excessive exposure in industrial situations or following acute poisoning, whereas organ tissue (lung, liver, kidney) cadmium concentrations may be useful in fatalities resulting from either acute or chronic poisoning. Cadmium concentrations in healthy persons without excessive cadmium exposure are generally less than 1 μg/L in either blood or urine. The ACGIH biological exposure indices for blood and urine cadmium levels are 5 μg/L and 5 μg/g creatinine, respectively, in random specimens. Persons who have sustained kidney damage due to chronic cadmium exposure often have blood or urine cadmium levels in a range of 25-50 μg/L or 25-75 μg/g creatinine, respectively. These ranges are usually 1000-3000 μg/L and 100-400 μg/g, respectively, in survivors of acute poisoning and may be substantially higher in fatal cases.
Treatment
A person with cadmium poisoning should seek immediate medical attention, both for treatment and supportive care.
For a non-chronic ingestive exposure, emetics or gastric lavage shortly after exposure can beneficially decontaminate the gastrointestinal system. Activated charcoal remains unproven. The US CDC does not recommend chelation therapy, in part because chelation may accentuate kidney damage.
For long-term exposure, considerable evidence indicates that the traditional chelator EDTA can reduce a body's overall cadmium load. Co-administered antioxidants, including nephroprotective glutathione, appear to improve efficacy. For patients with extremely fragile kidneys, limited evidence suggests that sauna sweat may differentially excrete the metal.
Epidemiology
In a mass cadmium poisoning in Japan, a marked prevalence for skeletal complications has been noted for older, postmenopausal women, however, the cause of the phenomenon is not fully understood, and is under investigation. Cadmium poisoning in postmenopausal women may result in an increased risk for osteoporosis. Current research has pointed to general malnourishment, as well as poor calcium metabolism relating to the women's age. Studies are pointing to damage of the mitochondria of kidney cells by cadmium as a key factor of the disease.
History
An experiment during the early 1960s involving the spraying of cadmium over Norwich was declassified in 2005 by the UK government, as documented in a BBC News article.
See also
Cobalt poisoning
Itai-itai disease
Citations
General and cited references
Shannon M. "Heavy Metal Poisoning", in Haddad LM, Shannon M, Winchester JF (editors): Clinical Management of Poisoning and Drug Overdose, Third Edition, 1998.
External links
ATSDR Case Studies in Environmental Medicine: Cadmium Toxicity U.S. Department of Health and Human Services
CDC - Cadmium - NIOSH Workplace Safety and Health Topic U.S. Department of Health and Human Services
National Pollutant Inventory - Cadmium and compounds
http://www.canoshweb.org/odp/html/cadmium.htm
After ‘Cadmium Rice,’ now ‘Lead’ and ‘Arsenic Rice’, New York Times
Poisoning
Contaminated farmland
Toxic effects of metals | Cadmium poisoning | [
"Chemistry",
"Environmental_science"
] | 2,431 | [
"Contaminated farmland",
"Water pollution"
] |
1,603,731 | https://en.wikipedia.org/wiki/Labdanum | Labdanum, also called ladanum, ladan, or ladanon, is a sticky brown resin obtained from the shrubs Cistus ladanifer (western Mediterranean) and Cistus creticus (eastern Mediterranean), species of rockrose. It was historically used in herbal medicine and is still used in the preparation of some perfumes and vermouths.
History
In ancient times, labdanum was collected by combing the beards and thighs of goats and sheep that had grazed on the cistus shrubs. Wooden instruments used were referred to in 19th-century Crete as ergastiri; a lambadistrion ("labdanum-gatherer") was a kind of rake to which a double row of leathern thongs were fixed instead of teeth. These were used to sweep the shrubs and collect the resin which was later extracted. It was collected by the shepherds and sold to coastal traders. The resin was used as an ingredient for incense, and medicinally to treat colds, coughs, menstrual problems and rheumatism.
Labdanum was produced on the banks of the Mediterranean in antiquity. The Book of Genesis contains two mentions of labdanum being carried to Egypt from Canaan. The word lot (לט "resin") in these two passages is usually interpreted as referring to labdanum on the basis of Semitic cognates.
Percy Newberry, a specialist on ancient Egypt, speculated that the false beard worn by Osiris and pharaohs may have originally represented a "labdanum-laden goat's beard". He also argued that the scepter of Osiris, which is usually interpreted as either a flail or a flabellum, was more likely an instrument for collecting labdanum similar to that used in nineteenth-century Crete.
Some scholars, such as Samuel Bochart, H.J. Abrahams, and Rabbi Saʻadiah ben Yosef Gaon (Saadya), 882–942, state that the mysterious שחלת (onycha), an ingredient in the holy incense (ketoret) mentioned in the Torah (Exodus 30: 34), was actually labdanum.
Modern uses
Labdanum is produced today mainly for the perfume industry. The raw resin is usually extracted by boiling the leaves and twigs. An absolute is also obtained by solvent extraction. An essential oil is produced by steam distillation. The raw gum is a black or sometimes dark brown, fragrant mass containing up to 20% or more of water. It is plastic but not pourable, and becomes brittle with age. The absolute is dark amber-green and very thick at room temperature. The fragrance is more refined than the raw resin. The odour is very rich, complex and tenacious. Labdanum is much valued in perfumery because of its resemblance to ambergris, which has been banned from use in many countries because it originates from the sperm whale, which is an endangered species. Labdanum is the main ingredient used when making the scent of amber in perfumery. Labdanum's odour is variously described as amber, sweet, woody, powdery, fruity, animalic, ambergris, dry musk, or leathery.
See also
Labdane
References
Resins
Perfume ingredients
Incense material | Labdanum | [
"Physics"
] | 689 | [
"Resins",
"Unsolved problems in physics",
"Incense material",
"Materials",
"Amorphous solids",
"Matter"
] |
1,604,816 | https://en.wikipedia.org/wiki/Steric%20factor | The steric factor, usually denoted ρ, is a quantity used in collision theory.
Also called the probability factor, the steric factor is defined as the ratio between the experimental value of the rate constant and the one predicted by collision theory. It can also be defined as the ratio between the pre-exponential factor and the collision frequency, and it is most often less than unity. Physically, the steric factor can be interpreted as the ratio of the cross section for reactive collisions to the total collision cross section.
Usually, the more complex the reactant molecules, the lower the steric factors. Nevertheless, some reactions exhibit steric factors greater than unity: the harpoon reactions, which involve atoms that exchange electrons, producing ions. The deviation from unity can have different causes: the molecules are not spherical, so different geometries are possible; not all the kinetic energy is delivered into the right spot; the presence of a solvent (when applied to solutions); and so on.
When collision theory is applied to reactions in solution, the solvent cage has an effect on the reactant molecules, as several collisions can take place in a single encounter, which leads to predicted preexponential factors being too large. ρ values greater than unity can be attributed to favorable entropic contributions.
Usually there is no simple way to accurately estimate steric factors without performing trajectory or scattering calculations. It is also more commonly known as the frequency factor.
Notes
Chemical kinetics
Physical chemistry | Steric factor | [
"Physics",
"Chemistry"
] | 296 | [
"Chemical reaction engineering",
"Applied and interdisciplinary physics",
"nan",
"Chemical kinetics",
"Physical chemistry",
"Physical chemistry stubs"
] |
1,605,201 | https://en.wikipedia.org/wiki/Solubility%20pump | In oceanic biogeochemistry, the solubility pump is a physico-chemical process that transports carbon as dissolved inorganic carbon (DIC) from the ocean's surface to its interior.
Overview
The solubility pump is driven by the coincidence of two processes in the ocean :
The solubility of carbon dioxide is a strong inverse function of seawater temperature (i.e. solubility is greater in cooler water)
The thermohaline circulation is driven by the formation of deep water at high latitudes where seawater is usually cooler and denser
Since deep water (that is, seawater in the ocean's interior) is formed under the same surface conditions that promote carbon dioxide solubility, it contains a higher concentration of dissolved inorganic carbon than might be expected from average surface concentrations. Consequently, these two processes act together to pump carbon from the atmosphere into the ocean's interior.
One consequence of this is that when deep water upwells in warmer, equatorial latitudes, it strongly outgasses carbon dioxide to the atmosphere because of the reduced solubility of the gas.
The solubility pump has a biological counterpart known as the biological pump. For an overview of both pumps, see Raven & Falkowski (1999).
Carbon dioxide solubility
Carbon dioxide, like other gases, is soluble in water. However, unlike many other gases (oxygen for instance), it reacts with water and forms a balance of several ionic and non-ionic species (collectively known as dissolved inorganic carbon, or DIC). These are dissolved free carbon dioxide (CO2 (aq)), carbonic acid (H2CO3), bicarbonate (HCO3−) and carbonate (CO32−), and they interact with water as follows :
The balance of these carbonate species (which ultimately affects the solubility of carbon dioxide), is dependent on factors such as pH, as shown in a Bjerrum plot. In seawater this is regulated by the charge balance of a number of positive (e.g. Na+, K+, Mg2+, Ca2+) and negative (e.g. CO32− itself, Cl−, SO42−, Br−) ions. Normally, the balance of these species leaves a net positive charge. With respect to the carbonate system, this excess positive charge shifts the balance of carbonate species towards negative ions to compensate. The result of which is a reduced concentration of the free carbon dioxide and carbonic acid species, which in turn leads to an oceanic uptake of carbon dioxide from the atmosphere to restore balance. Thus, the greater the positive charge imbalance, the greater the solubility of carbon dioxide. In carbonate chemistry terms, this imbalance is referred to as alkalinity.
In terms of measurement, four basic parameters are of key importance: Total inorganic carbon (TIC, T or CT), Total alkalinity (TALK or AT), pH, and pCO2. Measuring any two of these parameters allows for the determination of a wide range of pH-dependent species (including the above-mentioned species). This balance can be changed by a number of processes. For example, the air-sea flux of CO2, the dissolution/precipitation of CaCO3, or biological activity such as photosynthesis/respiration. Each of these has different effects on each of the four basic parameters, and together they exert strong influences on global cycles. The net and local charge of the oceans remains neutral during any chemical process.
Anthropogenic changes
The combustion of fossil fuels, land-use changes, and the production of cement have led to a flux of CO2 to the atmosphere. Presently, about one third (approximately 2 gigatons of carbon per year) of anthropogenic emissions of CO2 are believed to be entering the ocean. The solubility pump is the primary mechanism driving this flux, with the consequence that anthropogenic CO2 is reaching the ocean interior via high latitude sites of deep water formation (particularly the North Atlantic). Ultimately, most of the CO2 emitted by human activities will dissolve in the ocean, however the rate at which the ocean will take it up in the future is less certain.
In a study of carbon cycle up to the end of the 21st century, Cox et al. (2000) predicted that the rate of CO2 uptake will begin to saturate at a maximum rate at 5 gigatons of carbon per year by 2100. This was partially due to non-linearities in the seawater carbonate system, but also due to climate change. Ocean warming decreases the solubility of CO2 in seawater, slowing the ocean's response to emissions. Warming also acts to increase ocean stratification, isolating the surface ocean from deeper waters. Additionally, changes in the ocean's thermohaline circulation (specifically slowing) may act to decrease transport of dissolved CO2 into the deep ocean. However, the magnitude of these processes is still uncertain, preventing good long-term estimates of the fate of the solubility pump.
While ocean absorption of anthropogenic CO2 from the atmosphere acts to decrease climate change, it causes ocean acidification which is believed will have negative consequences for marine ecosystems.
See also
Alkalinity
Biological pump
Continental shelf pump
Ocean acidification
Thermohaline circulation
Total inorganic carbon
References
Aquatic ecology
Carbon cycle
Chemical oceanography
Geochemistry | Solubility pump | [
"Chemistry",
"Biology"
] | 1,117 | [
"Chemical oceanography",
"Aquatic ecology",
"Ecosystems",
"nan"
] |
1,605,292 | https://en.wikipedia.org/wiki/Data%20Transformation%20Services | Data Transformation Services (DTS) is a Microsoft database tool with a set of objects and utilities to allow the automation of extract, transform and load operations to or from a database. The objects are DTS packages and their components, and the utilities are called DTS tools. DTS was included with earlier versions of Microsoft SQL Server, and was almost always used with SQL Server databases, although it could be used independently with other databases.
DTS allows data to be transformed and loaded from heterogeneous sources using OLE DB, ODBC, or text-only files, into any supported database. DTS can also allow automation of data import or transformation on a scheduled basis, and can perform additional functions such as FTPing files and executing external programs. In addition, DTS provides an alternative method of version control and backup for packages when used in conjunction with a version control system, such as Microsoft Visual SourceSafe.
DTS has been superseded by SQL Server Integration Services in later releases of Microsoft SQL Server though there was some backwards compatibility and ability to run DTS packages in the new SSIS for a time.
History
In SQL Server versions 6.5 and earlier, database administrators (DBAs) used SQL Server Transfer Manager and Bulk Copy Program, included with SQL Server, to transfer data. These tools had significant shortcomings, and many DBAs used third-party tools such as Pervasive Data Integrator to transfer data more flexibly and easily. With the release of SQL Server 7 in 1998, "Data Transformation Services" was packaged with it to replace all these tools. The concept, design, and implementation of the Data Transformation Services was led by Stewart P. MacLeod (SQL Server Development Group Program Manager), Vij Rajarajan (SQL Server Lead Developer), and Ted Hart (SQL Server Lead Developer). The goal was to make it easier to import, export, and transform heterogeneous data and simplify the creation of data warehouses from operational data sources.
SQL Server 2000 expanded DTS functionality in several ways. It introduced new types of tasks, including the ability to FTP files, move databases or database components, and add messages into Microsoft Message Queue. DTS packages can be saved as a Visual Basic file in SQL Server 2000, and this can be expanded to save into any COM-compliant language. Microsoft also integrated packages into Windows 2000 security and made DTS tools more user-friendly; tasks can accept input and output parameters.
DTS comes with all editions of SQL Server 7 and 2000, but was superseded by SQL Server Integration Services in the Microsoft SQL Server 2005 release in 2005.
DTS packages
The DTS package is the fundamental logical component of DTS; every DTS object is a child component of the package. Packages are used whenever one modifies data using DTS. All the metadata about the data transformation is contained within the package. Packages can be saved directly in a SQL Server, or can be saved in the Microsoft Repository or in COM files. SQL Server 2000 also allows a programmer to save packages in a Visual Basic or other language file (when stored to a VB file, the package is actually scripted—that is, a VB script is executed to dynamically create the package objects and its component objects).
A package can contain any number of connection objects, but does not have to contain any. These allow the package to read data from any OLE DB-compliant data source, and can be expanded to handle other sorts of data. The functionality of a package is organized into tasks and steps.
A DTS Task is a discrete set of functionalities executed as a single step in a DTS package. Each task defines a work item to be performed as part of the data movement and data transformation process or as a job to be executed.
Data Transformation Services supplies a number of tasks that are part of the DTS object model and that can be accessed graphically through the DTS Designer or accessed programmatically. These tasks, which can be configured individually, cover a wide variety of data copying, data transformation and notification situations. For example, the following types of tasks represent some actions that you can perform by using DTS: executing a single SQL statement, sending an email, and transferring a file with FTP.
A step within a DTS package describes the order in which tasks are run and the precedence constraints that describe what to do in the case damage or of failure. These steps can be executed sequentially or in parallel.
Packages can also contain global variables which can be used throughout the package. SQL Server 2000 allows input and output parameters for tasks, greatly expanding the usefulness of global variables. DTS packages can be edited, password protected, scheduled for execution, and retrieved by version.
DTS tools
DTS tools packaged with SQL Server include the DTS wizards, DTS Designer, and DTS Programming Interfaces.
DTS wizards
The DTS wizards can be used to perform simple or common DTS tasks. These include the Import/Export Wizard and the Copy of Database Wizard. They provide the simplest method of copying data between OLE DB data sources. There is a great deal of functionality that is not available by merely using a wizard. However, a package created with a wizard can be saved and later altered with one of the other DTS tools.
A Create Publishing Wizard is also available to schedule packages to run at certain times. This only works if SQL Server Agent is running; otherwise the package will be scheduled, but will not be executed.
DTS Designer
The DTS Designer is a graphical tool used to build complex DTS Packages with workflows and event-driven logic. DTS Designer can also be used to edit and customize DTS Packages created with the DTS wizard.
Each connection and task in DTS Designer is shown with a specific icon. These icons are joined with precedence constraints, which specify the order and requirements for tasks to be run. One task may run, for instance, only if another task succeeds (or fails). Other tasks may run concurrently.
The DTS Designer has been criticized for having unusual quirks and limitations, such as the inability to visually copy and paste multiple tasks at one time. Many of these shortcomings have been overcome in SQL Server Integration Services, DTS's successor.
DTS Query Designer
A graphical tool used to build queries in DTS.
DTS Run Utility
DTS Packages can be run from the command line using the DTSRUN Utility.
The utility is invoked using the following syntax:
dtsrun /S server_name[\instance_name]
{ {/[~]U user_name [/[~]P password]} | /E }
]
{
{/[~]N package_name }
| {/[~]G package_guid_string}
| {/[~]V package_version_guid_string}
}
[/[~]M package_password]
[/[~]F filename]
[/[~]R repository_database_name]
[/A global_variable_name:typeid=value]
[/L log_file_name]
[/W NT_event_log_completion_status]
[/Z] [/!X] [/!D] [/!Y] [/!C]
]
When passing in parameters which are mapped to Global Variables, you are required to include the typeid. This is rather difficult to find on the Microsoft site. Below are the TypeIds used in passing in these values.
See also
OLAP
Data warehouse
Data mining
SQL Server Integration Services
Meta Data Services
References
External links
Microsoft SQL Server: Data Transformation Services (DTS)
SQL DTS unique DTS information resource
Understanding Microsoft Repository
DTS Videos & Training
Data Strategy
Microsoft database software
Data management
Extract, transform, load tools
Microsoft server technology | Data Transformation Services | [
"Technology"
] | 1,595 | [
"Data management",
"Data"
] |
5,642,583 | https://en.wikipedia.org/wiki/Lorenz%20system | The Lorenz system is a system of ordinary differential equations first studied by mathematician and meteorologist Edward Lorenz. It is notable for having chaotic solutions for certain parameter values and initial conditions. In particular, the Lorenz attractor is a set of chaotic solutions of the Lorenz system. The term "butterfly effect" in popular media may stem from the real-world implications of the Lorenz attractor, namely that tiny changes in initial conditions evolve to completely different trajectories. This underscores that chaotic systems can be completely deterministic and yet still be inherently impractical or even impossible to predict over longer periods of time. For example, even the small flap of a butterfly's wings could set the earth's atmosphere on a vastly different trajectory, in which for example a hurricane occurs where it otherwise would have not (see Saddle points). The shape of the Lorenz attractor itself, when plotted in phase space, may also be seen to resemble a butterfly.
Overview
In 1963, Edward Lorenz, with the help of Ellen Fetter who was responsible for the numerical simulations and figures, and Margaret Hamilton who helped in the initial, numerical computations leading up to the findings of the Lorenz model, developed a simplified mathematical model for atmospheric convection. The model is a system of three ordinary differential equations now known as the Lorenz equations:
The equations relate the properties of a two-dimensional fluid layer uniformly warmed from below and cooled from above. In particular, the equations describe the rate of change of three quantities with respect to time: is proportional to the rate of convection, to the horizontal temperature variation, and to the vertical temperature variation. The constants , , and are system parameters proportional to the Prandtl number, Rayleigh number, and certain physical dimensions of the layer itself.
The Lorenz equations can arise in simplified models for lasers, dynamos, thermosyphons, brushless DC motors, electric circuits, chemical reactions and forward osmosis. Interestingly, the same Lorenz equations were also derived in 1963 by Sauermann and Haken for a single-mode laser. In 1975, Haken realized that their equations derived in 1963 were mathematically equivalent to the original Lorenz equations. Haken's paper thus started a new field called laser chaos or optical chaos. The Lorenz equations are often called Lorenz-Haken equations in optical literature. Later on, it was also shown the complex version of Lorenz equations also had laser equivalent ones.
The Lorenz equations are also the governing equations in Fourier space for the Malkus waterwheel. The Malkus waterwheel exhibits chaotic motion where instead of spinning in one direction at a constant speed, its rotation will speed up, slow down, stop, change directions, and oscillate back and forth between combinations of such behaviors in an unpredictable manner.
From a technical standpoint, the Lorenz system is nonlinear, aperiodic, three-dimensional and deterministic. The Lorenz equations have been the subject of hundreds of research articles, and at least one book-length study.
Analysis
One normally assumes that the parameters , , and are positive. Lorenz used the values , , and . The system exhibits chaotic behavior for these (and nearby) values.
If then there is only one equilibrium point, which is at the origin. This point corresponds to no convection. All orbits converge to the origin, which is a global attractor, when .
A pitchfork bifurcation occurs at , and for two additional critical points appear at
These correspond to steady convection. This pair of equilibrium points is stable only if
which can hold only for positive if . At the critical value, both equilibrium points lose stability through a subcritical Hopf bifurcation.
When , , and , the Lorenz system has chaotic solutions (but not all solutions are chaotic). Almost all initial points will tend to an invariant setthe Lorenz attractora strange attractor, a fractal, and a self-excited attractor with respect to all three equilibria. Its Hausdorff dimension is estimated from above by the Lyapunov dimension (Kaplan-Yorke dimension) as , and the correlation dimension is estimated to be . The exact Lyapunov dimension formula of the global attractor can be found analytically under classical restrictions on the parameters:
The Lorenz attractor is difficult to analyze, but the action of the differential equation on the attractor is described by a fairly simple geometric model. Proving that this is indeed the case is the fourteenth problem on the list of Smale's problems. This problem was the first one to be resolved, by Warwick Tucker in 2002.
For other values of , the system displays knotted periodic orbits. For example, with it becomes a torus knot.
Connection to tent map
In Figure 4 of his paper, Lorenz plotted the relative maximum value in the z direction achieved by the system against the previous relative maximum in the direction. This procedure later became known as a Lorenz map (not to be confused with a Poincaré plot, which plots the intersections of a trajectory with a prescribed surface). The resulting plot has a shape very similar to the tent map. Lorenz also found that when the maximum value is above a certain cut-off, the system will switch to the next lobe. Combining this with the chaos known to be exhibited by the tent map, he showed that the system switches between the two lobes chaotically.
A Generalized Lorenz System
Over the past several years, a series of papers regarding high-dimensional Lorenz models have yielded a generalized Lorenz model, which can be simplified into the classical Lorenz model for three state variables or the following five-dimensional Lorenz model for five state variables:
A choice of the parameter has been applied to be consistent with the choice of the other parameters. See details in.
Simulations
Julia simulation
using Plots
# define the Lorenz attractor
@kwdef mutable struct Lorenz
dt::Float64 = 0.02
σ::Float64 = 10
ρ::Float64 = 28
β::Float64 = 8/3
x::Float64 = 2
y::Float64 = 1
z::Float64 = 1
end
function step!(l::Lorenz)
dx = l.σ * (l.y - l.x)
dy = l.x * (l.ρ - l.z) - l.y
dz = l.x * l.y - l.β * l.z
l.x += l.dt * dx
l.y += l.dt * dy
l.z += l.dt * dz
end
attractor = Lorenz()
# initialize a 3D plot with 1 empty series
plt = plot3d(
1,
xlim = (-30, 30),
ylim = (-30, 30),
zlim = (0, 60),
title = "Lorenz Attractor",
marker = 2,
)
# build an animated gif by pushing new points to the plot, saving every 10th frame
@gif for i=1:1500
step!(attractor)
push!(plt, attractor.x, attractor.y, attractor.z)
end every 10
Maple simulation
deq := [diff(x(t), t) = 10*(y(t) - x(t)), diff(y(t), t) = 28*x(t) - y(t) - x(t)*z(t), diff(z(t), t) = x(t)*y(t) - 8/3*z(t)]:
with(DEtools):
DEplot3d(deq, {x(t), y(t), z(t)}, t = 0 .. 100, [[x(0) = 10, y(0) = 10, z(0) = 10]], stepsize = 0.01, x = -20 .. 20, y = -25 .. 25, z = 0 .. 50, linecolour = sin(t*Pi/3), thickness = 1, orientation = [-40, 80], title = `Lorenz Chaotic Attractor`);
Maxima simulation
[sigma, rho, beta]: [10, 28, 8/3]$
eq: [sigma*(y-x), x*(rho-z)-y, x*y-beta*z]$
sol: rk(eq, [x, y, z], [1, 0, 0], [t, 0, 50, 1/100])$
len: length(sol)$
x: makelist(sol[k][2], k, len)$
y: makelist(sol[k][3], k, len)$
z: makelist(sol[k][4], k, len)$
draw3d(points_joined=true, point_type=-1, points(x, y, z), proportional_axes=xyz)$
MATLAB simulation
% Solve over time interval [0,100] with initial conditions [1,1,1]
% ''f'' is set of differential equations
% ''a'' is array containing x, y, and z variables
% ''t'' is time variable
sigma = 10;
beta = 8/3;
rho = 28;
f = @(t,a) [-sigma*a(1) + sigma*a(2); rho*a(1) - a(2) - a(1)*a(3); -beta*a(3) + a(1)*a(2)];
[t,a] = ode45(f,[0 100],[1 1 1]); % Runge-Kutta 4th/5th order ODE solver
plot3(a(:,1),a(:,2),a(:,3))
Mathematica simulation
Standard way:
tend = 50;
eq = {x'[t] == σ (y[t] - x[t]),
y'[t] == x[t] (ρ - z[t]) - y[t],
z'[t] == x[t] y[t] - β z[t]};
init = {x[0] == 10, y[0] == 10, z[0] == 10};
pars = {σ->10, ρ->28, β->8/3};
{xs, ys, zs} =
NDSolveValue[{eq /. pars, init}, {x, y, z}, {t, 0, tend}];
ParametricPlot3D[{xs[t], ys[t], zs[t]}, {t, 0, tend}]
Less verbose:
lorenz = NonlinearStateSpaceModel[{{σ (y - x), x (ρ - z) - y, x y - β z}, {}}, {x, y, z}, {σ, ρ, β}];
soln[t_] = StateResponse[{lorenz, {10, 10, 10}}, {10, 28, 8/3}, {t, 0, 50}];
ParametricPlot3D[soln[t], {t, 0, 50}]
Python simulation
import matplotlib.pyplot as plt
import numpy as np
def lorenz(xyz, *, s=10, r=28, b=2.667):
"""
Parameters
----------
xyz : array-like, shape (3,)
Point of interest in three-dimensional space.
s, r, b : float
Parameters defining the Lorenz attractor.
Returns
-------
xyz_dot : array, shape (3,)
Values of the Lorenz attractor's partial derivatives at *xyz*.
"""
x, y, z = xyz
x_dot = s*(y - x)
y_dot = r*x - y - x*z
z_dot = x*y - b*z
return np.array([x_dot, y_dot, z_dot])
dt = 0.01
num_steps = 10000
xyzs = np.empty((num_steps + 1, 3)) # Need one more for the initial values
xyzs[0] = (0., 1., 1.05) # Set initial values
# Step through "time", calculating the partial derivatives at the current point
# and using them to estimate the next point
for i in range(num_steps):
xyzs[i + 1] = xyzs[i] + lorenz(xyzs[i]) * dt
# Plot
ax = plt.figure().add_subplot(projection='3d')
ax.plot(*xyzs.T, lw=0.6)
ax.set_xlabel("X Axis")
ax.set_ylabel("Y Axis")
ax.set_zlabel("Z Axis")
ax.set_title("Lorenz Attractor")
plt.show()
R simulation
library(deSolve)
library(plotly)
# parameters
prm <- list(sigma = 10, rho = 28, beta = 8/3)
# initial values
varini <- c(
X = 1,
Y = 1,
Z = 1
)
Lorenz <- function (t, vars, prm) {
with(as.list(vars), {
dX <- prm$sigma*(Y - X)
dY <- X*(prm$rho - Z) - Y
dZ <- X*Y - prm$beta*Z
return(list(c(dX, dY, dZ)))
})
}
times <- seq(from = 0, to = 100, by = 0.01)
# call ode solver
out <- ode(y = varini, times = times, func = Lorenz,
parms = prm)
# to assign color to points
gfill <- function (repArr, long) {
rep(repArr, ceiling(long/length(repArr)))[1:long]
}
dout <- as.data.frame(out)
dout$color <- gfill(rainbow(10), nrow(dout))
# Graphics production with Plotly:
plot_ly(
data=dout, x = ~X, y = ~Y, z = ~Z,
type = 'scatter3d', mode = 'lines',
opacity = 1, line = list(width = 6, color = ~color, reverscale = FALSE)
)
Applications
Model for atmospheric convection
As shown in Lorenz's original paper, the Lorenz system is a reduced version of a larger system studied earlier by Barry Saltzman. The Lorenz equations are derived from the Oberbeck–Boussinesq approximation to the equations describing fluid circulation in a shallow layer of fluid, heated uniformly from below and cooled uniformly from above. This fluid circulation is known as Rayleigh–Bénard convection. The fluid is assumed to circulate in two dimensions (vertical and horizontal) with periodic rectangular boundary conditions.
The partial differential equations modeling the system's stream function and temperature are subjected to a spectral Galerkin approximation: the hydrodynamic fields are expanded in Fourier series, which are then severely truncated to a single term for the stream function and two terms for the temperature. This reduces the model equations to a set of three coupled, nonlinear ordinary differential equations. A detailed derivation may be found, for example, in nonlinear dynamics texts from , Appendix C; , Appendix D; or Shen (2016), Supplementary Materials.
Model for the nature of chaos and order in the atmosphere
The scientific community accepts that the chaotic features found in low-dimensional Lorenz models could represent features of the Earth's atmosphere (), yielding the statement of “weather is chaotic.” By comparison, based on the concept of attractor coexistence within the generalized Lorenz model and the original Lorenz model (), Shen and his co-authors proposed a revised view that “weather possesses both chaos and order with distinct predictability”. The revised view, which is a build-up of the conventional view, is used to suggest that “the chaotic and regular features found in theoretical Lorenz models could better represent features of the Earth's atmosphere”.
Resolution of Smale's 14th problem
Smale's 14th problem says, 'Do the properties of the Lorenz attractor exhibit that of a strange attractor?'. The problem was answered affirmatively by Warwick Tucker in 2002. To prove this result, Tucker used rigorous numerics methods like interval arithmetic and normal forms. First, Tucker defined a cross section that is cut transversely by the flow trajectories. From this, one can define the first-return map , which assigns to each the point where the trajectory of first intersects .
Then the proof is split in three main points that are proved and imply the existence of a strange attractor. The three points are:
There exists a region invariant under the first-return map, meaning .
The return map admits a forward invariant cone field.
Vectors inside this invariant cone field are uniformly expanded by the derivative of the return map.
To prove the first point, we notice that the cross section is cut by two arcs formed by . Tucker covers the location of these two arcs by small rectangles , the union of these rectangles gives . Now, the goal is to prove that for all points in , the flow will bring back the points in , in . To do that, we take a plan below at a distance small, then by taking the center of and using Euler integration method, one can estimate where the flow will bring in which gives us a new point . Then, one can estimate where the points in will be mapped in using Taylor expansion, this gives us a new rectangle centered on . Thus we know that all points in will be mapped in . The goal is to do this method recursively until the flow comes back to and we obtain a rectangle in such that we know that . The problem is that our estimation may become imprecise after several iterations, thus what Tucker does is to split into smaller rectangles and then apply the process recursively.
Another problem is that as we are applying this algorithm, the flow becomes more 'horizontal', leading to a dramatic increase in imprecision. To prevent this, the algorithm changes the orientation of the cross sections, becoming either horizontal or vertical.
Gallery
See also
Eden's conjecture on the Lyapunov dimension
Lorenz 96 model
List of chaotic maps
Takens' theorem
Notes
References
Shen, B.-W. (2015-12-21). "Nonlinear feedback in a six-dimensional Lorenz model: impact of an additional heating term". Nonlinear Processes in Geophysics. 22 (6): 749–764. doi:10.5194/npg-22-749-2015. ISSN 1607-7946.
Further reading
External links
Lorenz attractor by Rob Morris, Wolfram Demonstrations Project.
Lorenz equation on planetmath.org
Synchronized Chaos and Private Communications, with Kevin Cuomo. The implementation of Lorenz attractor in an electronic circuit.
Lorenz attractor interactive animation (you need the Adobe Shockwave plugin)
3D Attractors: Mac program to visualize and explore the Lorenz attractor in 3 dimensions
Lorenz Attractor implemented in analog electronic
Lorenz Attractor interactive animation (implemented in Ada with GTK+. Sources & executable)
Interactive web based Lorenz Attractor made with Iodide
Chaotic maps
Articles containing video clips
Articles with example Python (programming language) code
Articles with example MATLAB/Octave code
Articles with example Julia code | Lorenz system | [
"Mathematics"
] | 4,281 | [
"Functions and mappings",
"Mathematical objects",
"Mathematical relations",
"Chaotic maps",
"Dynamical systems"
] |
5,643,045 | https://en.wikipedia.org/wiki/Yitzhak%20Katznelson | Yitzhak Katznelson (; born 1934) is an Israeli mathematician.
Katznelson was born in Jerusalem. He received his doctoral degree from the University of Paris in 1956. He is a professor of mathematics at Stanford University.
He is the author of An Introduction to Harmonic Analysis, which won the Steele Prize for Mathematical Exposition in 2002.
In 2012 he became a fellow of the American Mathematical Society.
References
External links
An Introduction to Harmonic Analysis
1934 births
Living people
Israeli mathematicians
Jewish American scientists
Mathematical analysts
Stanford University Department of Mathematics faculty
Fellows of the American Mathematical Society
University of Paris alumni | Yitzhak Katznelson | [
"Mathematics"
] | 118 | [
"Mathematical analysis",
"Mathematical analysts"
] |
5,644,032 | https://en.wikipedia.org/wiki/Earth-centered%2C%20Earth-fixed%20coordinate%20system | The Earth-centered, Earth-fixed coordinate system (acronym ECEF), also known as the geocentric coordinate system, is a cartesian spatial reference system that represents locations in the vicinity of the Earth (including its surface, interior, atmosphere, and surrounding outer space) as X, Y, and Z measurements from its center of mass. Its most common use is in tracking the orbits of satellites and in satellite navigation systems for measuring locations on the surface of the Earth, but it is also used in applications such as tracking crustal motion.
The distance from a given point of interest to the center of Earth is called the geocentric distance, , which is a generalization of the geocentric radius, , not restricted to points on the reference ellipsoid surface.
The geocentric altitude is a type of altitude defined as the difference between the two aforementioned quantities: ; it is not to be confused for the geodetic altitude.
Conversions between ECEF and geodetic coordinates (latitude and longitude) are discussed at geographic coordinate conversion.
Structure
As with any spatial reference system, ECEF consists of an abstract coordinate system (in this case, a conventional three-dimensional right-handed system), and a geodetic datum that binds the coordinate system to actual locations on the Earth. The ECEF that is used for the Global Positioning System (GPS) is the geocentric WGS 84, which currently includes its own ellipsoid definition. Other local datums such as NAD 83 may also be used. Due to differences between datums, the ECEF coordinates for a location will be different for different datums, although the differences between most modern datums is relatively small, within a few meters.
The ECEF coordinate system has the following parameters:
The origin at the center of the chosen ellipsoid. In WGS 84, this is center of mass of the Earth.
The Z axis is the line between the North and South Poles, with positive values increasing northward. In WGS 84, this is the international reference pole (IRP), which does not exactly coincide with the Earth's rotational axis The slight "wobbling" of the rotational axis is known as polar motion, and can actually be measured against an ECEF.
The X axis is in the plane of the equator, passing through the origin and extending from 180° longitude (negative) to the prime meridian (positive); in WGS 84, this is the IERS Reference Meridian.
The Y axis is also in the plane of the equator, passing through extending from 90°W longitude (negative) to 90°E longitude (positive)
An example is the NGS data for a brass disk near Donner Summit, in California. Given the dimensions of the ellipsoid, the conversion from lat/lon/height-above-ellipsoid coordinates to X-Y-Z is straightforward—calculate the X-Y-Z for the given lat-lon on the surface of the ellipsoid and add the X-Y-Z vector that is perpendicular to the ellipsoid there and has length equal to the point's height above the ellipsoid. The reverse conversion is harder: given X-Y-Z can immediately get longitude, but no closed formula for latitude and height exists. See "Geodetic system." Using Bowring's formula in 1976 Survey Review the first iteration gives latitude correct within 10 degree as long as the point is within 10,000 meters above or 5,000 meters below the ellipsoid.
In astronomy
Geocentric coordinates can be used for locating astronomical objects in the Solar System in three dimensions along the Cartesian X, Y, and Z axes. They are differentiated from topocentric coordinates, which use the observer's location as the reference point for bearings in altitude and azimuth.
For nearby stars, astronomers use heliocentric coordinates, with the center of the Sun as the origin. The plane of reference can be aligned with the Earth's celestial equator, the ecliptic, or the Milky Way's galactic equator. These 3D celestial coordinate systems add actual distance as the Z axis to the equatorial, ecliptic, and galactic coordinate systems used in spherical astronomy.
See also
Earth-centered inertial (ECI)
Geodetic system
International Terrestrial Reference System and Frame (ITRS)
Orbital state vectors
Planetary coordinate system
References
External links
ECEF datum transformation Notes on converting ECEF coordinates to WGS-84 datum
Datum Transformations of GPS Positions Application Note Clearer notes on converting ECEF coordinates to WGS-84 datum
geodetic datum overview orientation of the coordinate system and additional information
GeographicLib includes a utility CartConvert which converts between geodetic and geocentric (ECEF) or local Cartesian (ENU) coordinates. This provides accurate results for all inputs including points close to the center of the Earth.
EPSG:4978
Global Positioning System
Astronomical coordinate systems | Earth-centered, Earth-fixed coordinate system | [
"Astronomy",
"Mathematics",
"Technology",
"Engineering"
] | 1,016 | [
"Wireless locating",
"Aerospace engineering",
"Astronomical coordinate systems",
"Aircraft instruments",
"Coordinate systems",
"Global Positioning System"
] |
5,644,561 | https://en.wikipedia.org/wiki/Marchenko%20equation | In mathematical physics, more specifically the one-dimensional inverse scattering problem, the Marchenko equation (or Gelfand-Levitan-Marchenko equation or GLM equation), named after Israel Gelfand, Boris Levitan and Vladimir Marchenko, is derived by computing the Fourier transform of the scattering relation:
Where is a symmetric kernel, such that which is computed from the scattering data. Solving the Marchenko equation, one obtains the kernel of the transformation operator from which the potential can be read off. This equation is derived from the Gelfand–Levitan integral equation, using the Povzner–Levitan representation.
Application to scattering theory
Suppose that for a potential for the Schrödinger operator , one has the scattering data , where are the reflection coefficients from continuous scattering, given as a function , and the real parameters are from the discrete bound spectrum.
Then defining
where the are non-zero constants, solving the GLM equation
for allows the potential to be recovered using the formula
See also
Lax pair
Notes
References
Eponymous equations of physics
Integral equations
Scattering theory | Marchenko equation | [
"Physics",
"Chemistry",
"Mathematics"
] | 219 | [
"Scattering theory",
"Equations of physics",
"Integral equations",
"Scattering stubs",
"Eponymous equations of physics",
"Mathematical objects",
"Equations",
"Scattering"
] |
13,578,878 | https://en.wikipedia.org/wiki/Impulse%20excitation%20technique | The impulse excitation technique (IET) is a non-destructive material characterization technique to determine the elastic properties and internal friction of a material of interest. It measures the resonant frequencies in order to calculate the Young's modulus, shear modulus, Poisson's ratio and internal friction of predefined shapes like rectangular bars, cylindrical rods and disc shaped samples. The measurements can be performed at room temperature or at elevated temperatures (up to 1700 °C) under different atmospheres.
The measurement principle is based on tapping the sample with a small projectile and recording the induced vibration signal with a piezoelectric sensor, microphone, laser vibrometer or accelerometer. To optimize the results a microphone or a laser vibrometer can be used as there is no contact between the test-piece and the sensor. Laser vibrometers are preferred to measure signals in vacuum. Afterwards, the acquired vibration signal in the time domain is converted to the frequency domain by a fast Fourier transformation. Dedicated software will determine the resonant frequency with high accuracy to calculate the elastic properties based on the classical beam theory.
Elastic properties
Different resonant frequencies can be excited dependent on the position of the support wires, the mechanical impulse and the microphone. The two most important resonant frequencies are the flexural which is controlled by the Young's modulus of the sample and the torsional which is controlled by the shear modulus for isotropic materials.
For predefined shapes like rectangular bars, discs, rods and grinding wheels, dedicated software calculates the sample's elastic properties using the sample dimensions, weight and resonant frequency (ASTM E1876-15).
Flexure mode
The first figure gives an example of a test-piece vibrating in the flexure
mode. This induced vibration is also referred as the out-of-plane vibration mode. The in-plane vibration will be excited by turning the sample 90° on the axis parallel to its length. The natural frequency of this flexural vibration mode is characteristic for the dynamic Young's modulus.
To minimize the damping of the test-piece, it has to be supported at the nodes where the vibration amplitude is zero. The test-piece is mechanically excited at one of the anti-nodes to cause maximum vibration.
Torsion mode
The second figure gives an example of a test-piece vibrating in the torsion mode. The natural frequency of this vibration is characteristic for the shear modulus.
To minimize the damping of the test-piece, it has to be supported at the center of both axis. The mechanical excitation has to be performed in one corner in order to twist the beam rather than flexing it.
Poisson's ratio
The Poisson's ratio is a measure in which a material tends to expand in directions perpendicular to the direction of compression. After measuring the Young's modulus and the shear modulus, dedicated software determines the Poisson's ratio using Hooke's law which can only be applied to isotropic materials according to the different standards.
Internal friction / Damping
Material damping or internal friction is characterized by the decay of the vibration amplitude of the sample in free vibration as the logarithmic decrement. The damping behaviour originates from anelastic processes occurring in a strained solid i.e. thermoelastic damping, magnetic damping, viscous damping, defect damping, ... For example, different materials defects (dislocations, vacancies, ...) can contribute to an increase in the internal friction between the vibrating defects and the neighboring regions.
Dynamic vs. static methods
Considering the importance of elastic properties for design and engineering applications, a number of experimental techniques are developed and these can be classified into 2 groups; static and dynamic methods. Statics methods (like the four-point bending test and nanoindentation) are based on direct measurements of stresses and strains during mechanical tests. Dynamic methods (like ultrasound spectroscopy and impulse excitation technique) provide an advantage over static methods because the measurements are relatively quick and simple and involve small elastic strains. Therefore, IET is very suitable for porous and brittle materials like ceramics and refractories. The technique can also be easily modified for high temperature experiments and only a small amount of material needs to be available.
Accuracy and uncertainty
The most important parameters to define the measurement uncertainty are the mass and dimensions of the sample. Therefore, each parameter has to be measured (and prepared) to a level of accuracy of 0.1%. Especially, the sample thickness is most critical (third power in the equation for Young's modulus). In that case, an overall accuracy of 1% can be obtained practically in most applications.
Applications
The impulse excitation technique can be used in a wide range of applications. Nowadays, IET equipment can perform measurements between −50 °C and 1700 °C in different atmospheres (air, inert, vacuum). IET is mostly used in research and as quality control tool to study the transitions as function of time and temperature.
A detailed insight into the material crystal structure can be obtained by studying the elastic and damping properties. For example, the interaction of dislocations and point defects in carbon steels are studied. Also the material damage accumulated during a thermal shock treatment can be determined for refractory materials. This can be an advantage in understanding the physical properties of certain materials.
Finally, the technique can be used to check the quality of systems. In this case, a reference piece is required to obtain a reference frequency spectrum. Engine blocks for example can be tested by tapping them and comparing the recorded signal with a pre-recorded signal of a reference engine block.
By using simple cluster analysis algorithms or principal component analysis, sample's pattern recognition is also achievable with a set of pre-recorded signals.
Experimental correlations
Rectangular bar
Young's modulus
with
E the Young's modulus
m the mass
ff the flexural frequency
b the width
L the length
t the thickness
T the correction factor
The correction factor can only be used if L/t ≥ 20!
Shear modulus
with
Note that we assume that b≥t
G the shear modulus
ft the torsional frequency
m the mass
b the width
L the length
t the thickness
R the correction factor
Cylindrical rod
Young's modulus
with
E the Young's modulus
m the mass
ff the flexural frequency
d the diameter
L the length
T the correction factorThe correction factor can only be used if L/d ≥ 20!Shear modulus
with
ft the torsional frequency
m the mass
d the diameter
L the length
Poisson ratio
If the Young's modulus and shear modulus are known, the Poisson's ratio can be calculated according to:
Damping coefficient
The induced vibration signal (in the time domain) is fitted as a sum of exponentially damped sinusoidal functions according to:
with
f the natural frequency
δ = kt the logarithmic decrement
In this case, the damping parameter Q−1 can be defined as:
with W the energy of the system
Extended IET applications: the Resonalyser Method
Isotropic versus orthotropic material behaviour
Isotropic elastic properties can be found by IET using the above described empirical formulas for the Young's modulus E, the shear modulus G and Poisson's ratio v. For isotropic materials the relation between strains and stresses in any point of flat sheets is given by the flexibility matrix [S] in the following expression:
In this expression, ε1 and ε2 are normal strains in the 1- and 2-direction and Υ12 is the shear strain. σ1 and σ2 are the normal stresses and τ12 is the shear stress. The orientation of the axes 1 and 2 in the above figure is arbitrary. This means that the values for E, G and v are the same in any material direction.
More complex material behaviour like orthotropic material behaviour can be identified by extended IET procedures. A material is called orthotropic when the elastic properties are symmetric with respect to a rectangular Cartesian system of axes. In case of a two dimensional state of stress, like in thin sheets, the stress-strain relations for orthotropic material become:
E1 and E2 are the Young's moduli in the 1- and 2-direction and G12 is the in-plane shear modulus. v12 is the major Poisson's ratio and v21 is the minor Poisson's ratio. The flexibility matrix [S] is symmetric. The minor Poisson's ratio can hence be found if E1, E2 and v12 are known.
The figure above shows some examples of common orthotropic materials: layered uni-directionally reinforced composites with fiber directions parallel to the plate edges, layered bi-directionally reinforced composites, short fiber reinforced composites with preference directions (like wooden particle boards), plastics with preference orientation, rolled metal sheets, and much more...
Extended IET for orthotropic material behaviour
Standard methods for the identification of the two Young's moduli E1 and E2 require two tensile, bending of IET tests, one on a beam cut along the 1-direction and one on a beam cut along the 2-direction. Major and minor Poisson's ratios can be identified if also the transverse strains are measured during the tensile tests. The identification of the in-plane shear modulus requires an additional in plane shearing test.
The "Resonalyser procedure" is an extension of the IET using an inverse method (also called "Mixed numerical experimental method"). The non destructive Resonalyser procedure allows a fast and accurate simultaneous identification of the 4 Engineering constants E1, E2, G12 and v12 for orthotropic materials. For the identification of the four orthotropic material constants, the first three natural frequencies of a rectangular test plate with constant thickness and the first natural frequency of two test beams with rectangular cross section must be measured. One test beam is cut along the longitudinal direction 1, the other one cut along the transversal direction 2 (see Figure on the right).
The Young's modulus of the test beams can be found using the bending IET formula for test beams with a rectangular cross section.
The ratio Width/Length of the test plate must be cut according to the following formula:
This ratio yields a so-called "Poisson plate". The interesting property of a Freely suspended Poisson plate is that the modal shapes that are associated with the 3 first resonance frequencies are fixed: the first resonance frequency is associated with a torsional modal shape, the second resonance frequency is associated with a saddle modal shape and the third resonance frequency is associated with a breathing modal shape.
364x364px
So, without the necessity to do an investigation to the nature of the modal shapes, the IET on a Poisson plate reveals the vibrational behaviour of a Poisson plate.
The question is now how to extract the orthotropic Engineering constants from the frequencies measured with IET on the beams and Poisson plate. This problem can be solved by an inverse method (also called" Mixed numerical/experimental method") based on a finite element (FE) computer model of the Poisson plate. A FE model allows computing resonance frequencies for a given set of material properties
In an inverse method, the material properties in the finite element model are updated in such a way that the computed resonance frequencies match the measured resonance frequencies.Problems with inverse methods are:
· The need of good starting values for the material properties
· Are the parameters converging to the correct physical solution?
· Is the solution unique?
The requirements to obtain good results are:
· The FE-model must be sufficiently accurate
· The IET measurements must be sufficiently accurate
· The starting values must be close enough to the final solution to avoid a local minimum (instead of a global minimum)
· The computed frequencies in the FE model of the Poisson plate must be sensitive for variations of all the material parameters
In the case the Young's moduli (obtained by IET) are fixed (as non variable parameters) in the inverse method procedure and if only the Poisson's ratio v12 and the in-plane shear modulus G12 are taken as variable parameters in the FE-model, the Resonalyser procedure''' satisfies all above requirements.
Indeed,
IET yields very accurate resonance frequencies, even with non-expert equipment,
a FE of a plate can be made very accurate by selecting a sufficiently fine element grid,
the knowledge of the modal shapes of a Poisson plate can be used to generate very good starting values using a virtual field method
and the first 3 natural frequencies of a Poisson plate are sensitive for variations of all the orthotropic Engineering constants.
Standards
ASTM E1876 - 15 Standard Test Method for Dynamic Youngs Modulus, Shear Modulus, and Poissons Ratio by Impulse Excitation of Vibration. www.astm.org.
ISO 12680-1:2005 - Methods of test for refractory products -- Part 1: Determination of dynamic Young's modulus (MOE) by impulse excitation of vibration. ISO.
DIN EN 843-2:2007 Advanced technical ceramics - Mechanical properties of monolithic ceramics at room temperature". webstore.ansi.org''.
References
Nondestructive testing
Quality control
Materials science
Continuum mechanics | Impulse excitation technique | [
"Physics",
"Materials_science",
"Engineering"
] | 2,804 | [
"Applied and interdisciplinary physics",
"Continuum mechanics",
"Classical mechanics",
"Materials science",
"Nondestructive testing",
"Materials testing",
"nan"
] |
13,584,590 | https://en.wikipedia.org/wiki/Electrokinetic%20phenomena | Electrokinetic phenomena are a family of several different effects that occur in heterogeneous fluids, or in porous bodies filled with fluid, or in a fast flow over a flat surface. The term heterogeneous here means a fluid containing particles. Particles can be solid, liquid or gas bubbles with sizes on the scale of a micrometer or nanometer. There is a common source of all these effects—the so-called interfacial 'double layer' of charges. Influence of an external force on the diffuse layer generates tangential motion of a fluid with respect to an adjacent charged surface. This force might be electric, pressure gradient, concentration gradient, or gravity. In addition, the moving phase might be either continuous fluid or dispersed phase.
Family
Various combinations of the driving force and moving phase determine various electrokinetic effects. According to J.Lyklema, the complete family of electrokinetic phenomena includes:
electrophoresis, as motion of charged particles under influence of electric field;
electro-osmosis, as motion of liquid in porous body under influence of electric field;
diffusiophoresis, as motion of particles under influence of a chemical potential gradient;
capillary osmosis, as motion of liquid in porous body under influence of the chemical potential gradient;
sedimentation potential, as electric field generated by sedimenting colloid particles;
streaming potential/current, as either electric potential or current generated by fluid moving through porous body, or relative to flat surface;
colloid vibration current, as electric current generated by particles moving in fluid under influence of ultrasound;
electric sonic amplitude, as ultrasound generated by colloidal particles in oscillating electric field.
Further reading
There are detailed descriptions of electrokinetic phenomena in many books on interface and colloid science.
See also
Isotachophoresis
Onsager reciprocal relations
Surface charge
Cationization of cotton
References
Colloidal chemistry
Condensed matter physics
Soft matter
Non-equilibrium thermodynamics
Electrochemistry | Electrokinetic phenomena | [
"Physics",
"Chemistry",
"Materials_science",
"Mathematics",
"Engineering"
] | 411 | [
"Colloidal chemistry",
"Non-equilibrium thermodynamics",
"Soft matter",
"Phases of matter",
"Materials science",
"Colloids",
"Surface science",
"Electrochemistry",
"Condensed matter physics",
"Matter",
"Dynamical systems"
] |
9,553,738 | https://en.wikipedia.org/wiki/Ensemble%20Kalman%20filter | The ensemble Kalman filter (EnKF) is a recursive filter suitable for problems with a large number of variables, such as discretizations of partial differential equations in geophysical models. The EnKF originated as a version of the Kalman filter for large problems (essentially, the covariance matrix is replaced by the sample covariance), and it is now an important data assimilation component of ensemble forecasting. EnKF is related to the particle filter (in this context, a particle is the same thing as an ensemble member) but the EnKF makes the assumption that all probability distributions involved are Gaussian; when it is applicable, it is much more efficient than the particle filter.
Introduction
The ensemble Kalman filter (EnKF) is a Monte Carlo implementation of the Bayesian update problem: given a probability density function (PDF) of the state of the modeled system (the prior, called often the forecast in geosciences) and the data likelihood, Bayes' theorem is used to obtain the PDF after the data likelihood has been taken into account (the posterior, often called the analysis). This is called a Bayesian update. The Bayesian update is combined with advancing the model in time, incorporating new data from time to time. The original Kalman filter, introduced in 1960, assumes that all PDFs are Gaussian (the Gaussian assumption) and provides algebraic formulas for the change of the mean and the covariance matrix by the Bayesian update, as well as a formula for advancing the mean and covariance in time provided the system is linear. However, maintaining the covariance matrix is not feasible computationally for high-dimensional systems. For this reason, EnKFs were developed. EnKFs represent the distribution of the system state using a collection of state vectors, called an ensemble, and replace the covariance matrix by the sample covariance computed from the ensemble. The ensemble is operated with as if it were a random sample, but the ensemble members are really not independent, as they all share the EnKF. One advantage of EnKFs is that advancing the PDF in time is achieved by simply advancing each member of the ensemble.
Derivation
Kalman filter
Let denote the -dimensional state vector of a model, and assume that it has Gaussian probability distribution with mean and covariance , i.e., its PDF is
Here and below, means proportional; a PDF is always scaled so that its integral over the whole space is one. This , called the prior, was evolved in time by running the model and now is to be updated to account for new data. It is natural to assume that the error distribution of the data is known; data have to come with an error estimate, otherwise they are meaningless. Here, the data is assumed to have Gaussian PDF with covariance and mean , where is the so-called observation matrix. The covariance matrix describes the estimate of the error of the data; if the random errors in the entries of the data vector are independent, is diagonal and its diagonal entries are the squares of the standard deviation (“error size”) of the error of the corresponding entries of the data vector . The value is what the value of the data would be for the state in the absence of data errors. Then the probability density of the data conditional of the system state , called the data likelihood, is
The PDF of the state and the data likelihood are combined to give the new probability density of the system state conditional on the value of the data (the posterior) by the Bayes theorem,
The data is fixed once it is received, so denote the posterior state by instead of and the posterior PDF by . It can be shown by algebraic manipulations that the posterior PDF is also Gaussian,
with the posterior mean and covariance given by the Kalman update formulas
where
is the so-called Kalman gain matrix.
Ensemble Kalman Filter
The EnKF is a Monte Carlo approximation of the Kalman filter, which avoids evolving the covariance matrix of the PDF of the state vector . Instead, the PDF is represented by an ensemble
is an matrix whose columns are the ensemble members, and it is called the prior ensemble. Ideally, ensemble members would form a sample from the prior distribution. However, the ensemble members are not in general independent except in the initial ensemble, since every EnKF step ties them together. They are deemed to be approximately independent, and all calculations proceed as if they actually were independent.
Replicate the data into an matrix
so that each column consists of the data vector plus a random vector from the -dimensional normal distribution . If, in addition, the columns of are a sample from the prior probability distribution, then the columns of
form a sample from the posterior probability distribution. To see this in the scalar case with : Let , and Then
.
The first sum is the posterior mean, and the second sum, in view of the independence, has a variance
,
which is the posterior variance.
The EnKF is now obtained simply by replacing the state covariance in Kalman gain matrix by the sample covariance computed from the ensemble members (called the ensemble covariance), that is:
Implementation
Basic formulation
Here we follow. Suppose the ensemble matrix and the data matrix are as above. The ensemble mean and the covariance are
where
and denotes the matrix of all ones of the indicated size.
The posterior ensemble is then given by
where the perturbed data matrix is as above.
Note that since is a covariance matrix, it is always positive semidefinite and usually positive definite, so the inverse above exists and the formula can be implemented by the Cholesky decomposition. In, is replaced by the sample covariance where and the inverse is replaced by a pseudoinverse, computed using the singular-value decomposition (SVD) .
Since these formulas are matrix operations with dominant Level 3 operations, they are suitable for efficient implementation using software packages such as LAPACK (on serial and shared memory computers) and ScaLAPACK (on distributed memory computers). Instead of computing the inverse of a matrix and multiplying by it, it is much better (several times cheaper and also more accurate) to compute the Cholesky decomposition of the matrix and treat the multiplication by the inverse as solution of a linear system with many simultaneous right-hand sides.
Observation matrix-free implementation
Since we have replaced the covariance matrix with ensemble covariance, this leads to a simpler formula where ensemble observations are directly used without explicitly specifying the matrix . More specifically, define a function of the form
The function is called the observation function or, in the inverse problems context, the forward operator. The value of is what the value of the data would be for the state assuming the measurement is exact. Then the posterior ensemble can be rewritten as
where
and
with
Consequently, the ensemble update can be computed by evaluating the observation function on each ensemble member once and the matrix does not need to be known explicitly. This formula holds also for an observation function with a fixed offset , which also does not need to be known explicitly. The above formula has been commonly used for a nonlinear observation function , such as the position of a hurricane vortex. In that case, the observation function is essentially approximated by a linear function from its values at ensemble members.
Implementation for a large number of data points
For a large number of data points, the multiplication by becomes a bottleneck. The following alternative formula is advantageous when the number of data points is large (such as when assimilating gridded or pixel data) and the data error covariance matrix is diagonal (which is the case when the data errors are uncorrelated), or cheap to decompose (such as banded due to limited covariance distance). Using the Sherman–Morrison–Woodbury formula
with
gives
which requires only the solution of systems with the matrix (assumed to be cheap) and of a system of size with right-hand sides. See for operation counts.
Further extensions
The EnKF version described here involves randomization of data. For filters without randomization of data, see.
Since the ensemble covariance is rank deficient (there are many more state variables, typically millions, than the ensemble members, typically less than a hundred), it has large terms for pairs of points that are spatially distant. Since in reality the values of physical fields at distant locations are not that much correlated, the covariance matrix is tapered off artificially based on the distance, which gives rise to localized EnKF algorithms. These methods modify the covariance matrix used in the computations and, consequently, the posterior ensemble is no longer made only of linear combinations of the prior ensemble.
For nonlinear problems, EnKF can create posterior ensemble with non-physical states. This can be alleviated by regularization, such as penalization of states with large spatial gradients.
For problems with coherent features, such as hurricanes, thunderstorms, firelines, squall lines, and rain fronts, there is a need to adjust the numerical model state by deforming the state in space (its grid) as well as by correcting the state amplitudes additively. In 2007, Ravela et al. introduce the joint position-amplitude adjustment model using ensembles, and systematically derive a sequential approximation which can be applied to both EnKF and other formulations. Their method does not make the assumption that amplitudes and position errors are independent or jointly Gaussian, as others do. The morphing EnKF employs intermediate states, obtained by techniques borrowed from image registration and morphing, instead of linear combinations of states.
Formally, EnKFs rely on the Gaussian assumption. In practice they can also be used for nonlinear problems, where the Gaussian assumption may not be satisfied. Related filters attempting to relax the Gaussian assumption in EnKF while preserving its advantages include filters that fit the state PDF with multiple Gaussian kernels, filters that approximate the state PDF by Gaussian mixtures, a variant of the particle filter with computation of particle weights by density estimation, and a variant of the particle filter with thick tailed data PDF to alleviate particle filter degeneracy.
See also
Data assimilation
Particle filter
Recursive Bayesian estimation
References
External links
EnKF webpage
TOPAZ, real-time forecasting of the North Atlantic ocean and Arctic sea-ice with the EnKF
EnKF-C, a compact framework for data assimilation into large-scale layered geophysical models with the EnKF
PDAF – Parallel Data Assimilation Framework – an open-source software for data assimilation providing different variants of the EnKF
Linear filters
Nonlinear filters
Bayesian statistics
Signal estimation
Monte Carlo methods | Ensemble Kalman filter | [
"Physics"
] | 2,204 | [
"Monte Carlo methods",
"Computational physics"
] |
9,553,854 | https://en.wikipedia.org/wiki/Artin%E2%80%93Rees%20lemma | In mathematics, the Artin–Rees lemma is a basic result about modules over a Noetherian ring, along with results such as the Hilbert basis theorem. It was proved in the 1950s in independent works by the mathematicians Emil Artin and David Rees; a special case was known to Oscar Zariski prior to their work.
An intuitive characterization of the lemma involves the notion that a submodule N of a module M over some ring A with specified ideal I holds a priori two topologies: one induced by the topology on M, and the other when considered with the I-adic topology over A. Then Artin-Rees dictates that these topologies actually coincide, at least when A is Noetherian and M finitely-generated.
One consequence of the lemma is the Krull intersection theorem. The result is also used to prove the exactness property of completion. The lemma also plays a key role in the study of ℓ-adic sheaves.
Statement
Let I be an ideal in a Noetherian ring R; let M be a finitely generated R-module and let N a submodule of M. Then there exists an integer k ≥ 1 so that, for n ≥ k,
Proof
The lemma immediately follows from the fact that R is Noetherian once necessary notions and notations are set up.
For any ring R and an ideal I in R, we set (B for blow-up.) We say a decreasing sequence of submodules is an I-filtration if ; moreover, it is stable if for sufficiently large n. If M is given an I-filtration, we set ; it is a graded module over .
Now, let M be a R-module with the I-filtration by finitely generated R-modules. We make an observation
is a finitely generated module over if and only if the filtration is I-stable.
Indeed, if the filtration is I-stable, then is generated by the first terms and those terms are finitely generated; thus, is finitely generated. Conversely, if it is finitely generated, say, by some homogeneous elements in , then, for , each f in can be written as
with the generators in . That is, .
We can now prove the lemma, assuming R is Noetherian. Let . Then are an I-stable filtration. Thus, by the observation, is finitely generated over . But is a Noetherian ring since R is. (The ring is called the Rees algebra.) Thus, is a Noetherian module and any submodule is finitely generated over ; in particular, is finitely generated when N is given the induced filtration; i.e., . Then the induced filtration is I-stable again by the observation.
Krull's intersection theorem
Besides the use in completion of a ring, a typical application of the lemma is the proof of the Krull's intersection theorem, which says: for a proper ideal I in a commutative Noetherian ring that is either a local ring or an integral domain. By the lemma applied to the intersection , we find k such that for ,
Taking , this means or . Thus, if A is local, by Nakayama's lemma. If A is an integral domain, then one uses the determinant trick (that is a variant of the Cayley–Hamilton theorem and yields Nakayama's lemma):
In the setup here, take u to be the identity operator on N; that will yield a nonzero element x in A such that , which implies , as is a nonzerodivisor.
For both a local ring and an integral domain, the "Noetherian" cannot be dropped from the assumption: for the local ring case, see local ring#Commutative case. For the integral domain case, take to be the ring of algebraic integers (i.e., the integral closure of in ). If is a prime ideal of A, then we have: for every integer . Indeed, if , then for some complex number . Now, is integral over ; thus in and then in , proving the claim.
Footnotes
References
gives a somehow more precise version of the Artin–Rees lemma.
External links
Commutative algebra
Lemmas in algebra
Module theory
Theorems in ring theory | Artin–Rees lemma | [
"Mathematics"
] | 917 | [
"Theorems in algebra",
"Lemmas in algebra",
"Fields of abstract algebra",
"Module theory",
"Commutative algebra",
"Lemmas"
] |
9,555,814 | https://en.wikipedia.org/wiki/Trans-Earth%20injection | A trans-Earth injection (TEI) is a propulsion maneuver used to set a spacecraft on a trajectory which will intersect the Earth's sphere of influence, usually putting the spacecraft on a free return trajectory.
The maneuver is performed by a rocket engine.
From the Moon
The spacecraft is usually in a parking orbit around the Moon at the time of TEI, in which case the burn is timed so that its midpoint is opposite the Earth's location upon arrival. Uncrewed space probes have also performed this maneuver from the Moon starting with Luna 16's direct ascent traverse from the lunar surface in 1970.
On the Apollo missions, it was performed by the restartable Service Propulsion System (SPS) engine on the Service Module after the undocking of the (LM) Lunar Module if provided. An Apollo TEI burn lasted approximately 150 seconds, providing a posigrade velocity increase of 1,000 m/s (3,300 ft/s). It was first performed by the Apollo 8 mission on December 25, 1968. It was last performed by the propulsion module of Chandrayaan-3 mission during 13 October 2023
List of missions that performed a Trans-Earth injection
Total 17 missions have performed such a maneuver. NASA has performed it the most (10 times), followed by Soviet Union (3 times), China (3 times), and India (once). These missions are in order,
Apollo 8
Apollo 10
Apollo 11
Apollo 12
Luna 16
Apollo 14
Apollo 15
Luna 20
Apollo 16
Apollo 17
Luna 24
Clementine
Chang'e 5-T1
Chang'e 5
Artemis 1
Chandrayaan-3
Chang'e 6
From outside the Earth-Moon system
In 2004, from outside the Earth-Moon system, the Stardust probe comet dust return mission performed TEI after visiting Comet Wild 2.
See also
Lunar orbit insertion
Trans-lunar injection
Trans-Mars injection
References
Astrodynamics
Spacecraft propulsion
Exploration of the Moon
Apollo program | Trans-Earth injection | [
"Astronomy",
"Engineering"
] | 397 | [
"Aerospace engineering",
"Astrodynamics",
"Astronomy stubs",
"Spacecraft stubs"
] |
9,556,852 | https://en.wikipedia.org/wiki/Quark%20epoch | In physical cosmology, the quark epoch was the period in the evolution of the early universe when the fundamental interactions of gravitation, electromagnetism, the strong interaction and the weak interaction had taken their present forms, but the temperature of the universe was still too high to allow quarks to bind together to form hadrons. The quark epoch began approximately 10−12 seconds after the Big Bang, when the preceding electroweak epoch ended as the electroweak interaction separated into the weak interaction and electromagnetism. During the quark epoch, the universe was filled with a dense, hot quark–gluon plasma, containing quarks, leptons and their antiparticles. Collisions between particles were too energetic to allow quarks to combine into mesons or baryons. The quark epoch ended when the universe was about 10−6 seconds old, when the average energy of particle interactions had fallen below the binding energy of hadrons. The following period, when quarks became confined within hadrons, is known as the hadron epoch.
See also
Timeline of the early universe
Chronology of the universe
Cosmology
References
Further reading
Big Bang
Physical cosmology | Quark epoch | [
"Physics",
"Astronomy"
] | 249 | [
"Cosmogony",
"Astronomical sub-disciplines",
"Big Bang",
"Theoretical physics",
"Astrophysics",
"Physical cosmology"
] |
17,538,047 | https://en.wikipedia.org/wiki/Willam%E2%80%93Warnke%20yield%20criterion | The Willam–Warnke yield criterion is a function that is used to predict when failure will occur in concrete and other cohesive-frictional materials such as rock, soil, and ceramics. This yield criterion has the functional form
where is the first invariant of the Cauchy stress tensor, and are the second and third invariants of the deviatoric part of the Cauchy stress tensor. There are three material parameters ( - the uniaxial compressive strength, – the uniaxial tensile strength, - the equibiaxial compressive strength) that have to be determined before the Willam-Warnke yield criterion may be applied to predict failure.
In terms of , the Willam-Warnke yield criterion can be expressed as
where is a function that depends on and the three material parameters and depends only on the material parameters. The function can be interpreted as the friction angle which depends on the Lode angle (). The quantity is interpreted as a cohesion pressure. The Willam-Warnke yield criterion may therefore be viewed as a combination of the Mohr–Coulomb and the Drucker–Prager yield criteria.
Willam-Warnke yield function
In the original paper, the three-parameter Willam-Warnke yield function was expressed as
where is the first invariant of the stress tensor, is the second invariant of the deviatoric part of the stress tensor, is the yield stress in uniaxial compression, and is the Lode angle given by
The locus of the boundary of the stress surface in the deviatoric stress plane is expressed in polar coordinates by the quantity which is given by
where
The quantities and describe the position vectors at the locations and can be expressed in terms of as (here is the failure stress under equi-biaxial compression and is the failure stress under uniaxial tension)
The parameter in the model is given by
The Haigh-Westergaard representation of the Willam-Warnke yield condition can be
written as
where
Modified forms of the Willam-Warnke yield criterion
An alternative form of the Willam-Warnke yield criterion in Haigh-Westergaard coordinates is the Ulm-Coussy-Bazant form:
where
and
The quantities are interpreted as friction coefficients. For the yield surface to be convex, the Willam-Warnke yield criterion requires that and .
See also
Yield (engineering)
Yield surface
Plasticity (physics)
References
Chen, W. F. (1982). Plasticity in Reinforced Concrete. McGraw Hill. New York.
External links
Kaspar Willam and E.P. Warnke (1974). Constitutive model for the triaxial behavior of concrete
Palko, J. L. (1993). Interactive reliability model for whisker-toughened ceramics
The ‘‘Chunnel’’ Fire. I: Chemoplastic softening in rapidly heated concrete by Franz-Josef Ulm, Olivier Coussy, and Zdeneˇk P. Bazˇant.
Plasticity (physics)
Solid mechanics
Yield criteria | Willam–Warnke yield criterion | [
"Materials_science"
] | 624 | [
"Deformation (mechanics)",
"Plasticity (physics)"
] |
17,543,136 | https://en.wikipedia.org/wiki/Structured%20Stream%20Transport | In computer networking, Structured Stream Transport (SST) is an experimental transport protocol that provides an ordered, reliable byte stream abstraction similar to TCP's, but enhances and optimizes stream management to permit applications to use streams in a much more fine-grained fashion than is feasible with TCP streams.
External links
SST home page
Transport layer protocols
Network protocols | Structured Stream Transport | [
"Technology"
] | 77 | [
"Computing stubs",
"Computer network stubs"
] |
17,545,070 | https://en.wikipedia.org/wiki/List%20of%20instruments%20used%20in%20microbiological%20sterilization%20and%20disinfection | This is a list of instruments used in microbiological sterilization and disinfection.
Instrument list
References
Microbiology equipment | List of instruments used in microbiological sterilization and disinfection | [
"Biology"
] | 28 | [
"Microbiology equipment"
] |
17,545,102 | https://en.wikipedia.org/wiki/Instruments%20used%20in%20medical%20laboratories | This is a list of instruments used in general in laboratories, including:
Biochemistry
Microbiology
Pharmacology
Instrument list
Image gallery
References
Medical equipment
Biochemistry methods
Laboratory equipment
Microbiology equipment
Clinical pathology | Instruments used in medical laboratories | [
"Chemistry",
"Biology"
] | 39 | [
"Biochemistry methods",
"Medical equipment",
"Microbiology equipment",
"Biochemistry",
"Medical technology"
] |
17,545,130 | https://en.wikipedia.org/wiki/Instruments%20used%20in%20pathology | Instruments used specially in pathology are as follows:
Instrument list
Gallery
References
Medical equipment
Pathology | Instruments used in pathology | [
"Biology"
] | 18 | [
"Medical equipment",
"Pathology",
"Medical technology"
] |
17,545,297 | https://en.wikipedia.org/wiki/Instruments%20used%20in%20obstetrics%20and%20gynecology | The following is a list of instruments that are used in modern obstetrics and gynaecology.
See also
Instruments used in general medicine
Medical procedure
Women's health
References
Obstetrics
Gynaecology
Medical devices
Medical lists | Instruments used in obstetrics and gynecology | [
"Biology"
] | 48 | [
"Medical devices",
"Medical technology"
] |
17,545,444 | https://en.wikipedia.org/wiki/Aflibercept | Aflibercept, sold under the brand names Eylea and Zaltrap among others, is a medication used to treat wet macular degeneration and metastatic colorectal cancer. It was developed by Regeneron Pharmaceuticals.
It is an inhibitor of vascular endothelial growth factor (VEGF). Aflibercept is a recombinant fusion protein consisting of the extracellular domains of human VEGF receptor 1 and 2 fused to the Fc portion of human IgG1. By acting as a soluble decoy for the natural VEGF receptors, aflibercept inhibits their activation, thereby reducing angiogenesis.
Medical uses
Aflibercept (Eylea) is indicated for the treatment of people with neovascular (wet) age-related macular degeneration, macular edema following retinal vein occlusion, diabetic macular edema, diabetic retinopathy, and retinopathy of prematurity.
Aflibercept (Zaltrap), in combination with fluorouracil, leucovorin, and irinotecan (known as FOLFIRI), is indicated for the treatment of people with metastatic colorectal cancer that is resistant to, or has progressed following, an oxaliplatin-containing regimen.
It is used for the treatment of wet macular degeneration and is administered as an intravitreal injection, that is, into the eye. For cancer treatment, it is given intravenously in combination with fluorouracil, leucovorin, and irinotecan.
In July 2014, aflibercept (Eylea) was approved for the treatment of people with visual impairment due to diabetic macular edema In May 2019, the US FDA expanded the indication for aflibercept to include all stages of diabetic retinopathy.
In February 2023, the US FDA approved aflibercept (Eylea) as a treatment for retinopathy of prematurity.
Contraindications
Aflibercept (Eylea) is contraindicated in people with infections or active inflammations of or near the eye, while aflibercept (Zaltrap) has no contraindications.
Adverse effects
Common adverse effects of the eye formulation include conjunctival hemorrhage, eye pain, cataract, vitreous detachment, floaters, and ocular hypertension.
Aflibercept (Zaltrap) has adverse effects typical of anti-cancer drugs, such as reduced blood cell count (leukopenia, neutropenia, thrombocytopenia), gastrointestinal disorders like diarrhea and abdominal pain, and fatigue. Another common effect is hypertension (increased blood pressure).
Interactions
No interactions are described for either formulation.
Mechanism of action
In wet macular degeneration, abnormal blood vessels grow in the choriocapillaris, a layer of capillaries in the eye, leading to blood and protein leakage below the macula.
Aflibercept (Zaltrap) binds to circulating VEGFs and acts like a "VEGF trap". It thereby inhibits the activity of the vascular endothelial growth factor subtypes VEGF-A and VEGF-B, as well as to placental growth factor (PGF), inhibiting the growth of new blood vessels in the choriocapillaris or the tumour, respectively. The aim of the cancer treatment, so to speak, is to starve the tumor.
Composition
Aflibercept is a recombinant fusion protein consisting of vascular endothelial growth factor (VEGF)-binding portions from the extracellular domains of human VEGF receptors 1 and 2, that are fused to the Fc portion of the human IgG1 immunoglobulin.
History
Regeneron commenced clinical testing of aflibercept in cancer in 2001. In 2003, Regeneron signed a major deal with Aventis to develop aflibercept in the field of cancer. In 2004 Regeneron started testing the compound, locally delivered, in proliferative eye diseases, and in 2006 Regeneron and Bayer signed an agreement to develop the eye indications.
Society and culture
Legal status
In November 2011, the US Food and Drug Administration (FDA) approved aflibercept for the treatment of wet macular degeneration.
In August 2012, the US FDA approved aflibercept (Zaltrap) for use in combination with 5-fluorouracil, folinic acid and irinotecan to treat adults with metastatic colorectal cancer that is resistant to, or has progressed following, an oxaliplatin‑containing regimen. To avoid confusion with the version that is injected into the eye, the FDA assigned a new name, ziv-aflibercept, to the active ingredient.
In November 2012, the European Medicines Agency (EMA) approved aflibercept (Eylea) for the treatment of wet macular degeneration.
In February 2013, the European Medicines Agency (EMA) approved aflibercept (Zaltrap) for the treatment of adults with metastatic colorectal cancer for whom treatment based on oxaliplatin has not worked or the cancer got worse. Aflibercept (Zaltrap) is used with irinotecan, 5-fluorouracil, and folinic acid.
In August 2023, the FDA approved aflibercept (Eylea) for the treatment of wet age-related macular degeneration, diabetic macular edema, and diabetic retinopathy.
Biosimilars
Yesafili was approved for medical use in the European Union in September 2023.
In May 2024, aflibercept-jbvf (Yesafili) and aflibercept-yszy (Opuviz) were approved for medical use in the United States.
Aflibercept-mrbb (Ahzantive) was approved for medical use in the United States in June 2024. It is a biosimilar to Eylea.
In August 2024, aflibercept-abzv (Enzeevu) was approved for medical use in the United States. It is a biosimilar to Eylea.
In August 2024, aflibercept-ayyh (Pavblu) was approved for medical use in the United States. It is a biosimilar to Eylea.
In September 2024, the Committee for Medicinal Products for Human Use (CHMP) of the European Medicines Agency adopted a positive opinion, recommending the granting of a marketing authorization for the medicinal product Opuviz, intended for the treatment of neovascular (wet) age-related macular degeneration, visual impairment due to macular edema secondary to retinal vein occlusion (branch RVO or central RVO), visual impairment due to diabetic macular edema (DME) and visual impairment due to myopic choroidal neovascularization (myopic CNV). The applicant for this medicinal product is Samsung Bioepis NL B.V. Opuviz is a biosimilar medicinal product that is highly similar to the reference product Eylea.
In September 2024, the CHMP adopted a positive opinion, recommending the granting of a marketing authorization for the medicinal product Afqlir, intended for the treatment of neovascular (wet) age-related macular degeneration, visual impairment due to macular edema secondary to retinal vein occlusion (branch RVO or central RVO), visual impairment due to diabetic macular edema (DME) and visual impairment due to myopic choroidal neovascularization (myopic CNV). The applicant for this medicinal product is Sandoz GmbH. Afqlir is a biosimilar medicinal product that is highly similar to the reference product Eylea. Afqlir was authorized for use in the EU in November 2024.
In November 2024, the CHMP adopted a positive opinion, recommending the granting of a marketing authorization for the medicinal product Ahzantive, intended for the treatment of neovascular (wet) age-related macular degeneration, visual impairment due to macular edema secondary to retinal vein occlusion (branch RVO or central RVO), visual impairment due to diabetic macular edema (DME) and visual impairment due to myopic choroidal neovascularisation (myopic CNV). The applicant for this medicinal product is Klinge Biopharma GmbH. Ahzantive is a biosimilar medicinal product that is highly similar to the reference product Eylea. Ahzantive was approved for medical use in the European Union in January 2025.
In November 2024, the CHMP adopted a positive opinion, recommending the granting of a marketing authorization for the medicinal product Baiama, intended for the treatment of neovascular (wet) age-related macular degeneration, visual impairment due to macular edema secondary to retinal vein occlusion (branch RVO or central RVO), visual impairment due to diabetic macular edema (DME) and visual impairment due to myopic choroidal neovascularisation (myopic CNV). The applicant for this medicinal product is Formycon AG. Baiama is a biosimilar medicinal product that is highly similar to the reference product Eylea. Baiama was approved for medical use in the European Union in January 2025.
In December 2024, the CHMP adopted a positive opinion, recommending the granting of a marketing authorization for the medicinal product Eydenzelt, intended for the treatment of neovascular (wet) age-related macular degeneration (AMD), visual impairment due to macular oedema secondary to retinal vein occlusion (branch RVO or central RVO), visual impairment due to diabetic macular oedema (DME) and visual impairment due to myopic choroidal neovascularisation (myopic CNV). The applicant for this medicinal product is Celltrion Healthcare Hungary Kft. Eydenzelt is a biosimilar medicinal product. It is highly similar to the reference product Eylea.
Economics
In March 2015, aflibercept was one of a group of drugs delisted from the UK Cancer Drugs Fund. In 2017, injections of aflibercept (HCPCS code J0178) were responsible for the most billing to Medicare Part B, at .
Research
In March 2011, aflibercept failed its primary endpoint of overall survival in the Vital phase III trial for second-line treatment of locally advanced or metastatic non-small cell lung cancer, although it improved the secondary endpoint of progression-free survival.
In April 2011, aflibercept improved its primary endpoint of overall survival in the Velour phase III clinical trial for second-line treatment for metastatic colorectal cancer.
Aflibercept was also in a phase III trial for hormone-refractory metastatic prostate cancer .
A 2016 Cochrane Review examined outcomes comparing aflibercept versus ranibizumab injections in over 2400 people with neovascular AMD, from two randomized controlled trials. Both treatment options yielded similar improvements in visual acuity and morphological outcomes, though the authors note that the aflibercept treatment regimen has the potential to reduce treatment burden and risks from frequent injections.
A 2017 review update studying the effects of anti-VEGF drugs on diabetic macular edema found that while all three studied treatments have advantages over laser therapy, there was moderate evidence that aflibercept is significantly favored in all measured efficacy outcomes over ranibizumab and bevacizumab, after one year, longer term advantages were unclear.
References
External links
Angiogenesis inhibitors
Drugs developed by Bayer
Engineered proteins
Ophthalmology drugs
Sanofi | Aflibercept | [
"Biology"
] | 2,553 | [
"Angiogenesis",
"Angiogenesis inhibitors"
] |
17,546,914 | https://en.wikipedia.org/wiki/SAM1 | SAM1, or "Semiempirical ab initio Model 1", is a semiempirical quantum chemistry method for computing molecular properties. It is an implementation the general Neglect of Differential Diatomic Overlap (NDDO) integral approximation, and is efficient and accurate. Related methods are AM1, PM3 and the older MNDO.
SAM1 was developed by M.J.S. Dewar and co-workers at the University of Texas and the University of Florida. Papers describing the implementation of the method and its results were published in 1993 and 1994. The method is implemented in the AMPAC program produced by Semichem
SAM1 builds on the success of the Dewar-style semiempirical models by adding two new aspects to the AM1/PM3 formalism:
Two-electron repulsion integrals (TERIs) are computed from a minimal basis set of contracted Gaussian functions, as opposed to the previously used multipole expansion. Note that the NDDO approximation is still in effect, and that only a few of the possible TERIs are explicitly computed. The values of the explicit TERIs are scaled using empirically-derived functions to obtain experimentally relevant results.
One-center two-electron repulsion integrals (OCTEs) are derived initially to reproduce atomic properties. These values are then fixed and carried forward as further elemental parameterization proceeds.
The performance of SAM1 for C, H, O, N, F, Cl, Br, and I was claimed to be superior to other semiempirical methods. Especially noteworthy were the smaller systematic errors for heats for formation.
.
See also
Semi-empirical quantum chemistry method
NDDO
References
Semiempirical quantum chemistry methods | SAM1 | [
"Chemistry"
] | 357 | [
"Quantum chemistry stubs",
"Quantum chemistry",
"Theoretical chemistry stubs",
"Computational chemistry",
"Physical chemistry stubs",
"Semiempirical quantum chemistry methods"
] |
12,023,401 | https://en.wikipedia.org/wiki/Dermaseptin | Dermaseptins are a family of peptides isolated from skin of the frog genus Phyllomedusa. The sequence of the dermaseptins varies greatly but due to the presence of lysine residues all are cationic and most have the potential to form amphipathic helices in water or when integrated with the lipid bilayer of the bacterial membrane. Clear separation of two lobes of positive and negative intramolecular electrostatic potential is thought to be important in cytotoxic activity. Dermaseptins are typically 27-34 amino acid residues in length and were the first vertebrate peptides demonstrated as having a lethal effect on the filamentous fungi implicated in severe opportunistic infections accompanying immunodeficiency syndrome and immunosuppressive drug therapy.
Dermaseptin use in a novel drug delivery system has been proposed. The system is based on the affinity of dermaseptins for the plasma membrane of human erythrocytes. After transient loading of the cells with the non-toxic dermaseptin S4 analogue K4–S4(1–13)a, the peptide is transported in the systemic circulation to distant microbial targets. Upon reaching a microorganism for which it has greater affinity the dermaseptin derivative is spontaneously transferred to the microbial membrane where it exerts its membrane-lytic activity.
See also
Antimicrobial resistance
References
Amphibian toxins
Peptides | Dermaseptin | [
"Chemistry"
] | 303 | [
"Biomolecules by chemical classification",
"Peptides",
"Molecular biology"
] |
12,026,279 | https://en.wikipedia.org/wiki/Pancreatic%20ribonuclease%20family | Pancreatic ribonuclease family (, RNase, RNase I, RNase A, pancreatic RNase, ribonuclease I, endoribonuclease I, ribonucleic phosphatase, alkaline ribonuclease, ribonuclease, gene S glycoproteins, Ceratitis capitata alkaline ribonuclease, SLSG glycoproteins, gene S locus-specific glycoproteins, S-genotype-assocd. glycoproteins, ribonucleate 3'-pyrimidino-oligonucleotidohydrolase) is a superfamily of pyrimidine-specific endonucleases found in high quantity in the pancreas of certain mammals and of some reptiles.
Specifically, the enzymes are involved in endonucleolytic cleavage of 3'-phosphomononucleotides and 3'-phosphooligonucleotides ending in C-P or U-P with 2',3'-cyclic phosphate intermediates. Ribonuclease can unwind the RNA helix by complexing with single-stranded RNA; the complex arises by an extended multi-site cation-anion interaction between lysine and arginine residues of the enzyme and phosphate groups of the nucleotides.
Notable family members
Bovine pancreatic ribonuclease is the best-studied member of the family and has served as a model system in work related to protein folding, disulfide bond formation, protein crystallography and spectroscopy, and protein dynamics. The human genome contains 8 genes that share the structure and function with bovine pancreatic ribonuclease, with 5 additional pseudo-genes. The structure and dynamics of these enzymes are related to their diverse biological functions.
Other proteins belonging to the pancreatic ribonuclease superfamily include: bovine seminal vesicle and brain ribonucleases; kidney non-secretory ribonucleases; liver-type ribonucleases; angiogenin, which induces vascularisation of normal and malignant tissues; eosinophil cationic protein, a cytotoxin and helminthotoxin with ribonuclease activity; and frog liver ribonuclease and frog sialic acid-binding lectin. The sequence of pancreatic ribonucleases contains four conserved disulfide bonds and three amino acid residues involved in the catalytic activity.
Human genes
Human genes encoding proteins containing this domain include:
ANG,
RNASE1, RNASE10, RNASE12, RNASE2, RNASE3, RNASE4, RNASE6, RNASE7, and RNASE8.
Cytotoxicity
Some members of the pancreatic ribonuclease family have cytotoxic effects. Mammalian cells are protected from these effects due to their extremely high affinity for ribonuclease inhibitor (RI), which protects cellular RNA from degradation by pancreatic ribonucleases. Pancreatic ribonucleases that are not inhibited by RI are approximately as toxic as alpha-sarcin, diphtheria toxin, or ricin.
Two pancreatic ribonucleases isolated from the oocytes of the Northern leopard frog - amphinase and ranpirnase - are not inhibited by RI and show differential cytotoxicity against tumor cells. Ranpirnase was studied in a Phase III clinical trial as a treatment candidate for mesothelioma, but the trial did not demonstrate statistical significance against primary endpoints.
References
Ribonucleases
EC 3.1.27
Protein domains | Pancreatic ribonuclease family | [
"Biology"
] | 795 | [
"Protein domains",
"Protein classification"
] |
14,708,063 | https://en.wikipedia.org/wiki/Enthalpy%E2%80%93entropy%20compensation | In thermodynamics, enthalpy–entropy compensation is a specific example of the compensation effect. The compensation effect refers to the behavior of a series of closely related chemical reactions (e.g., reactants in different solvents or reactants differing only in a single substituent), which exhibit a linear relationship between one of the following kinetic or thermodynamic parameters for describing the reactions:
Between the logarithm of the pre-exponential factors (or prefactors) and the activation energies where the series of closely related reactions are indicated by the index , are the preexponential factors, are the activation energies, is the gas constant, and , are constants.
Between enthalpies and entropies of activation (enthalpy–entropy compensation) where are the enthalpies of activation and are the entropies of activation.
Between the enthalpy and entropy changes of a series of similar reactions (enthalpy–entropy compensation) where are the enthalpy changes and are the entropy changes.
When the activation energy is varied in the first instance, we may observe a related change in pre-exponential factors. An increase in tends to compensate for an increase in , which is why we call this phenomenon a compensation effect. Similarly, for the second and third instances, in accordance with the Gibbs free energy equation, with which we derive the listed equations, scales proportionately with . The enthalpy and entropy compensate for each other because of their opposite algebraic signs in the Gibbs equation.
A correlation between enthalpy and entropy has been observed for a wide variety of reactions. The correlation is significant because, for linear free-energy relationships (LFERs) to hold, one of three conditions for the relationship between enthalpy and entropy for a series of reactions must be met, with the most common encountered scenario being that which describes enthalpy–entropy compensation. The empirical relations above were noticed by several investigators beginning in the 1920s, since which the compensatory effects they govern have been identified under different aliases.
Related terms
Many of the more popular terms used in discussing the compensation effect are specific to their field or phenomena. In these contexts, the unambiguous terms are preferred. The misapplication of and frequent crosstalk between fields on this matter has, however, often led to the use of inappropriate terms and a confusing picture. For the purposes of this entry different terms may refer to what may seem to be the same effect, but that either a term is being used as a shorthand (isokinetic and isoequilibrium relationships are different, yet are often grouped together synecdochically as isokinetic relationships for the sake of brevity) or is the correct term in context. This section should aid in resolving any uncertainties. (see Criticism section for more on the variety of terms)
compensation effect/rule : umbrella term for the observed linear relationship between: (i) the logarithm of the preexponential factors and the activation energies, (ii) enthalpies and entropies of activation, or (iii) between the enthalpy and entropy changes of a series of similar reactions.
enthalpy-entropy compensation : the linear relationship between either the enthalpies and entropies of activation or the enthalpy and entropy changes of a series of similar reactions.
isoequilibrium relation (IER), isoequilibrium effect : On a Van 't Hoff plot, there exists a common intersection point describing the thermodynamics of the reactions. At the isoequilibrium temperature , all the reactions in the series should have the same equilibrium constant ()
isokinetic relation (IKR), isokinetic effect : On an Arrhenius plot, there exists a common intersection point describing the kinetics of the reactions. At the isokinetic temperature , all the reactions in the series should have the same rate constant ()
isoequilibrium temperature : used for thermodynamic LFERs; refers to in the equations where it possesses dimensions of temperature
isokinetic temperature : used for kinetic LFERs; refers to in the equations where it possesses dimensions of temperature
kinetic compensation : an increase in the preexponential factors tends to compensate for the increase in activation energy:
Meyer–Neldel rule (MNR) : primarily used in materials science and condensed matter physics; the MNR is often stated as the plot of the logarithm of the preexponential factor against activation energy is linear:
where is the preexponential factor, is the activation energy, σ is the conductivity, and is the Boltzmann constant, and is temperature.
Mathematics
Enthalpy–entropy compensation as a requirement for LFERs
Linear free-energy relationships (LFERs) exist when the relative influence of changing substituents on one reactant is similar to the effect on another reactant, and include linear Hammett plots, Swain–Scott plots, and Brønsted plots. LFERs are not always found to hold, and to see when one can expect them to, we examine the relationship between the free-energy differences for the two reactions under comparison. The extent to which the free energy of the new reaction is changed, via a change in substituent, is proportional to the extent to which the reference reaction was changed by the same substitution. A ratio of the free-energy differences is the reaction quotient or constant .
The above equation may be rewritten as the difference () in free-energy changes ():
Substituting the Gibbs free-energy equation () into the equation above yields a form that makes clear the requirements for LFERs to hold.
One should expect LFERs to hold if one of three conditions are met:
's are coincidentally the same for both the new reaction under study and the reference reaction, and the 's are linearly proportional for the two reactions being compared.
's are coincidentally the same for both the new reaction under study and the reference reaction, and the 's are linearly proportional for the two reactions being compared.
's and 's are linearly related to each other for both the reference reaction and the new reaction.
The third condition describes the enthalpy–entropy effect and is the condition most commonly met.
Isokinetic and isoequilibrium temperature
For most reactions the activation enthalpy and activation entropy are unknown, but, if these parameters have been measured and a linear relationship is found to exist (meaning an LFER was found to hold), the following equation describes the relationship between and :
Inserting the Gibbs free-energy equation and combining like terms produces the following equation:
where is constant regardless of substituents and is different for each substituent.
In this form, has the dimension of temperature and is referred to as the isokinetic (or isoequilibrium) temperature.
Alternately, the isokinetic (or isoequilibrium) temperature may be reached by observing that, if a linear relationship is found, then the difference between the s for any closely related reactants will be related to the difference between 's for the same reactants:
Using the Gibbs free-energy equation,
In both forms, it is apparent that the difference in Gibbs free-energies of activations () will be zero when the temperature is at the isokinetic (or isoequilibrium) temperature and hence identical for all members of the reaction set at that temperature.
Beginning with the Arrhenius equation and assuming kinetic compensation (obeying ), the isokinetic temperature may also be given by
The reactions will have approximately the same value of their rate constant at an isokinetic temperature.
History
In a 1925 paper, F.H. Constable described the linear relationship observed for the reaction parameters of the catalytic dehydrogenation of primary alcohols with copper-chromium oxide.
Phenomenon explained
The foundations of the compensation effect are still not fully understood though many theories have been brought forward. Compensation of Arrhenius processes in solid-state materials and devices can be explained quite generally from the statistical physics of aggregating fundamental excitations from the thermal bath to surmount a barrier whose activation energy is significantly larger than the characteristic energy of the excitations used (e.g., optical phonons). To rationalize the occurrences of enthalpy-entropy compensation in protein folding and enzymatic reactions, a Carnot-cycle model in which a micro-phase transition plays a crucial role was proposed. In drug receptor binding, it has been suggested that enthalpy-entropy compensation arises due to an intrinsic property of hydrogen bonds. A mechanical basis for solvent-induced enthalpy-entropy compensation has been put forward and tested at the dilute gas limit. There is some evidence of enthalpy-entropy compensation in biochemical or metabolic networks particularly in the context of intermediate-free coupled reactions or processes. However, a single general statistical mechanical explanation applicable to all compensated processes has not yet been developed.
Criticism
Kinetic relations have been observed in many systems and, since their conception, have gone by many terms, among which are the Meyer-Neldel effect or rule, the Barclay-Butler rule, the theta rule, and the Smith-Topley effect. Generally, chemists will talk about the isokinetic relation (IKR), from the importance of the isokinetic (or isoequilibrium) temperature, condensed matter physicists and material scientists use the Meyer-Neldel rule, and biologists will use the compensation effect or rule.
An interesting homework problem appears following Chapter 7: Structure-Reactivity Relationships in Kenneth Connors's textbook Chemical Kinetics: The Study of Reaction Rates:
From the last four digits of the office telephone numbers of the faculty in your department, systematically construct pairs of "rate constants" as two-digit numbers times 10−5 s−1 at temperatures 300 K and 315 K (obviously the larger rate constant of each pair to be associated with the higher temperature). Make a two-point Arrhenius plot for each faculty member, evaluating and . Examine the plot of against for evidence of an isokinetic relationship.
The existence of any real compensation effect has been widely derided in recent years and attributed to the analysis of interdependent factors and chance. Because the physical roots remain to be fully understood, it has been called into question whether compensation is a truly physical phenomenon or a coincidence due to trivial mathematical connections between parameters. The compensation effect has been criticized in other respects, namely for being the result of random experimental and systematic errors producing the appearance of compensation. The principal complaint lodged states that compensation is an artifact of data from a limited temperature range or from a limited range for the free energies.
In response to the criticisms, investigators have stressed that compensatory phenomena are real, but appropriate and in-depth data analysis is always needed. The F-test has been used to such an aim, and it minimizes the deviations of points constrained to pass through an isokinetic temperature to the deviation of the points from the unconstrained line is achieved by comparing the mean deviations of points. Appropriate statistical tests should be performed as well. W. Linert wrote in a 1983 paper:
There are few topics in chemistry in which so many misunderstandings and controversies have arisen as in connection with the so-called isokinetic relationship (IKR) or compensation law. Up to date, a great many chemists appear to be inclined to dismiss the IKR as being accidental. The crucial problem is that the activation parameters are mutually dependent because of their determination from the experimental data. Therefore, it has been stressed repeatedly, the isokinetic plot (i.e., against ) is unfit in principle to substantiate a claim of an isokinetic relationship. At the same time, however, it is a fatal error to dismiss the IKR because of that fallacy.
Common among all defenders is the agreement that stringent criteria for the assignment of true compensation effects must be adhered to.
References
Thermodynamics
Chemical thermodynamics | Enthalpy–entropy compensation | [
"Physics",
"Chemistry",
"Mathematics"
] | 2,496 | [
"Chemical thermodynamics",
"Thermodynamics",
"Dynamical systems"
] |
14,709,342 | https://en.wikipedia.org/wiki/Cation%20channel%20superfamily | The transmembrane cation channel superfamily was defined in InterPro and Pfam as the family of tetrameric ion channels. These include the sodium, potassium, calcium, ryanodine receptor, HCN, CNG, CatSper, and TRP channels. This large group of ion channels apparently includes families , , , and of the TCDB transporter classification.
They are described as minimally having two transmembrane helices flanking a loop which determines the ion selectivity of the channel pore. Many eukaryotic channels have four additional transmembrane helices (TM) (), related to or vestigial of voltage gating. The proteins with only two transmembrane helices () are most commonly found in bacteria. This also includes the 2-TM inward-rectifier potassium channels () found primarily in eukaryotes. There are commonly additional regulatory domains which serve to regulate ion conduction and channel gating. The pores may also be homotetramers or heterotetramers; where heterotetramers may be encoded as distinct genes or as multiple pore domains within a single polypeptide. The HVCN1 and Putative tyrosine-protein phosphatase proteins do not contain an expected ion conduction pore domain, but rather have homology only to the voltage sensor domain of voltage gated ion channels.
Human channels with 6 TM helices
Cation
Transient receptor potential
Canonical
TRPC1; TRPC3; TRPC4; TRPC5; TRPC6; TRPC7
Melastatin
TRPM1; TRPM2; TRPM3; TRPM4; TRPM5; TRPM6; TRPM7; TRPM8
Vanilloid
TRPV1; TRPV2; TRPV3; TRPV4; TRPV5; TRPV6
Mucolipin
MCOLN1; MCOLN2; MCOLN3;
Ankyrin
TRPA1
TRPP
PKD1L3;
Calcium
Voltage-dependent
CACNA1A; CACNA1B; CACNA1C; CACNA1D; CACNA1E; CACNA1F; CACNA1G; CACNA1H; CACNA1I; CACNA1S
Sperm
CATSPER1; CATSPER2; CATSPER3; CATSPER4
Ryanodine receptor
RYR1; RYR2; RYR3
Potassium
Voltage-gated potassium
Delayed rectifier
Kvα1.x - Shaker-related: Kv1.1 (KCNA1), Kv1.2 (KCNA2), Kv1.3 (KCNA3), Kv1.5 (KCNA5), Kv1.6 (KCNA6), Kv1.7 (KCNA7), Kv1.8 (KCNA10)
Kvα2.x - Shab-related: Kv2.1 (KCNB1), Kv2.2 (KCNB2)
Kvα3.x - Shaw-related: Kv3.1 (KCNC1), Kv3.2 (KCNC2)
Kvα7.x: Kv7.1 (KCNQ1) - KvLQT1, Kv7.2 (KCNQ2), Kv7.3 (KCNQ3), Kv7.4 (KCNQ4), Kv7.5 (KCNQ5)
Kvα10.x: Kv10.1 (KCNH1)
A-type potassium
Kvα1.x - Shaker-related: Kv1.4 (KCNA4)
Kvα3.x - Shaw-related: Kv3.3 (KCNC3), Kv3.4 (KCNC4)
Kvα4.x - Shal-related: Kv4.1 (KCND1), Kv4.2 (KCND2), Kv4.3 (KCND3)
Outward-rectifying
Kvα10.x: Kv10.2 (KCNH5)
Inwardly-rectifying
Kvα11.x - ether-a-go-go potassium channels: Kv11.1 (KCNH2) - hERG, Kv11.2 (KCNH6), Kv11.3 (KCNH7)
Slowly activating
Kvα12.x: Kv12.1 (KCNH8), Kv12.2 (KCNH3), Kv12.3 (KCNH4)
Modifier/silencer
Kvα5.x: Kv5.1 (KCNF1)
Kvα6.x: Kv6.1 (KCNG1), Kv6.2 (KCNG2), Kv6.3 (KCNG3), Kv6.4 (KCNG4)
Kvα8.x: Kv8.1 (KCNV1), Kv8.2 (KCNV2)
Kvα9.x: Kv9.1 (KCNS1), Kv9.2 (KCNS2), Kv9.3 (KCNS3)
Calcium-activated
BK
KCa1.1 (BK, Slo1, Maxi-K, )
SK
KCa2.x: KCa2.1 (KCNN1) - SK1, KCa2.2 (KCNN2) - SK2, KCa2.3 (KCNN3) - SK3
KCa3.x: KCa3.1 (KCNN4) - SK4
KCa4.x: KCa4.1 (KCNT1) - SLACK, KCa4.2 (KCNT2) - SLICK
IK
KCa3.1 (IKCa1, SK4, )
Other subfamilies
KCa5.1 (Slo3, )
Inward-rectifier potassium
Sodium
NALCN
SCN1A; SCN2A; SCN2A2; SCN3A; SCN4A; SCN5A; SCN7A; SCN8A; SCN9A; SCN10A; SCN11A
SLC9A10; SLC9A11
Cyclic nucleotide-gated
CNGA1; CNGA2; CNGA3; CNGA4
CNGB1; CNGB3
HCN1; HCN2; HCN3; HCN4
ITPR1; ITPR2; ITPR3
Proton
HVCN1
Related proteins
TPTE, part of the larger Voltage sensitive phosphatase family
Human channels with 2 TM helices in each subunit
Potassium
Tandem pore domain potassium channel
KCNK1; KCNK2; KCNK3; KCNK4; KCNK5; KCNK6; KCNK7; KCNK9; KCNK10; KCNK12; KCNK13; KCNK15; KCNK16; KCNK17; KCNK18
Non-human channels
Two-pore
TPCN1
TPCN2
Pore-only potassium
KcsA
Ligand-gated potassium
GluR0
Voltage-gated potassium
KvAP
Prokaryotic KCa
Kch
MthK
TrkA/TrkH
KtrAB
GsuK
TM1088
Voltage and cyclic nucleotide gated potassium
MlotiK1
Sodium
NaChBac
NaVAb
NaVAe1
NaVAp
NaVMm
Non-selective
NaK
Prokaryotic inward-rectifier potassium
KirBac
Engineered
NaK2CNG
NaK2K
References
External links
Protein domains
Protein families
Transmembrane proteins
Ion channels | Cation channel superfamily | [
"Chemistry",
"Biology"
] | 1,606 | [
"Protein classification",
"Protein domains",
"Protein families",
"Neurochemistry",
"Ion channels"
] |
14,710,128 | https://en.wikipedia.org/wiki/Potassium%20channel%20tetramerisation%20domain | K+ channel tetramerisation domain is the N-terminal, cytoplasmic tetramerisation domain (T1) of voltage-gated K+ channels. It defines molecular determinants for subfamily-specific assembly of alpha-subunits into functional tetrameric channels. It is distantly related to the BTB/POZ domain .
Potassium channels
Potassium channels are the most diverse group of the ion channel family. They are important in shaping the action potential, and in neuronal excitability and plasticity. The potassium channel family is composed of several functionally distinct isoforms, which can be broadly separated into 2 groups: the practically non-inactivating 'delayed' group and the rapidly inactivating 'transient' group.
These are all highly similar proteins, with only small amino acid changes causing the diversity of the voltage-dependent gating mechanism, channel conductance and toxin binding properties. Each type of K+ channel is activated by different signals and conditions depending on their type of regulation: some open in response to depolarisation of the plasma membrane; others in response to hyperpolarisation or an increase in intracellular calcium concentration; some can be regulated by binding of a transmitter, together with intracellular kinases; while others are regulated by GTP-binding proteins or other second messengers. In eukaryotic cells, K+ channels are involved in neural signalling and generation of
the cardiac rhythm, act as effectors in signal transduction pathways involving G protein-coupled receptors (GPCRs)
and may have a role in target cell lysis by cytotoxic T-lymphocytes. In prokaryotic cells, they play a role in the maintenance of ionic homeostasis.
Alpha subunits of the channels
All K+ channels discovered so far possess a core of alpha subunits, each comprising either one or two copies of a highly conserved pore loop domain (P-domain). The P-domain contains the sequence (T/SxxTxGxG), which has
been termed the K+ selectivity sequence. In families that contain one P-domain, four subunits assemble to form a selective pathway for K+ across the membrane. However, it remains unclear how the 2 P-domain subunits assemble to form a selective pore. The functional diversity of these families can arise through homo- or hetero-associations of alpha subunits or association with auxiliary cytoplasmic beta subunits. K+ channel subunits containing one pore domain can be assigned into one of two superfamilies: those that possess six transmembrane (TM) domains and those that possess only two TM domains. The six TM domain superfamily can be further subdivided into conserved gene families: the voltage-gated (Kv) channels; the KCNQ channels (originally known as KvLQT channels); the EAG-like K+ channels; and three types of calcium (Ca)-activated K+ channels (BK, IK and SK). The 2TM domain family comprises inward-rectifying K+ channels. In addition, there are K+ channel alpha-subunits that possess two P-domains. These are usually highly regulated K+ selective leak channels.
The Kv family can be divided into several subfamilies on the basis of sequence similarity and function. Four of these subfamilies, Kv1 (Shaker), Kv2 (Shab), Kv3 (Shaw) and Kv4 (Shal), consist of pore-forming alpha subunits that associate with different types of beta subunit. Each alpha subunit comprises six hydrophobic TM domains with a P-domain between the fifth and sixth, which partially resides in the membrane. The fourth TM domain has positively charged residues at every third residue and acts as a voltage sensor, which triggers the conformational change that opens the channel pore in response to a displacement in membrane potential. More recently, 4 new electrically-silent alpha subunits have been
cloned: Kv5 (KCNF), Kv6 (KCNG), Kv8 and Kv9 (KCNS). These subunits do not themselves possess any functional activity, but appear to form heteromeric channels with Kv2 subunits, and thus modulate Shab channel activity. When highly expressed, they inhibit channel activity, but at lower levels show more specific modulatory actions.
Tetramerization domain
The N-terminal, cytoplasmic tetramerization domain (T1) of voltage-gated potassium channels encodes molecular determinants for subfamily-specific assembly of alpha-subunits into functional tetrameric channels. This domain is found in a subset of a larger group of proteins that contain the BTB/POZ domain.
Human proteins containing this domain
BTBD10; KCNA1; KCNA10; KCNA2; KCNA3; KCNA4; KCNA5; KCNA6;
KCNA7; KCNB1; KCNB2; KCNC1; KCNC2; KCNC3; KCNC4; KCND1;
KCND2; KCND3; KCNF1; KCNG1; KCNG2; KCNG3; KCNG4; KCNRG;
KCNS1; KCNS2; KCNS3; KCNV1; KCNV2; KCTD1; KCTD10; KCTD11;
KCTD12; KCTD13; KCTD14; KCTD15; KCTD16; KCTD17; KCTD18; KCTD19;
KCTD2; KCTD20; KCTD21; KCTD3; KCTD4; KCTD5; KCTD6; KCTD7;
KCTD8; KCTD9; SHKBP1; TNFAIP1;
References
Further reading
Protein domains
Transmembrane proteins | Potassium channel tetramerisation domain | [
"Biology"
] | 1,235 | [
"Protein domains",
"Protein classification"
] |
14,710,215 | https://en.wikipedia.org/wiki/GRAM%20domain | The GRAM domain is found in glucosyltransferases, myotubularins and other membrane-associated proteins. The structure of the GRAM domain is similar to that found in PH domains.
Proteins containing GRAM domains are found in all eukaryotes and bacteria, but not archaea. Various GRAM domains can bind proteins or lipids.
Human proteins containing this domain
GRAMD1A; GRAMD1B; GRAMD1C; GRAMD2A; GRAMD2B; GRAMD4; MTM1; MTMR1; MTMR2; NCOA7; NSMAF; OXR1; SBF1; SBF2; TBC1D8; TBC1D8B; TBC1D9; TBC1D9B; WBP2; WBP2NL; dJ439F8.1;
References
Protein domains
Protein families
Peripheral membrane proteins | GRAM domain | [
"Biology"
] | 193 | [
"Protein families",
"Protein domains",
"Protein classification"
] |
14,710,808 | https://en.wikipedia.org/wiki/Animal%20heme-dependent%20peroxidases | Animal heme-dependent peroxidases is a family of peroxidases. Peroxidases are found in bacteria, fungi, plants and animals. On the basis of sequence similarity, a number of animal heme peroxidases can be categorized as members of a superfamily: myeloperoxidase (MPO); eosinophil peroxidase (EPO); lactoperoxidase (LPO); thyroid peroxidase (TPO); prostaglandin H synthase (PGHS); and peroxidasin.
Function
Myeloperoxidase (MPO) plays a major role in the oxygen-dependent microbicidal system of neutrophils. EPO from eosinophilic granulocytes participates in immunological reactions, and potentiates tumor necrosis factor (TNF) production and hydrogen peroxide release by human monocyte-derived macrophages. MPO (and possibly EPO) primarily use Cl−ions and H2O2 to form hypochlorous acid (HOCl), which can effectively kill bacteria or parasites. In secreted fluids, LPO catalyses the oxidation of thiocyanate ions (SCN−) by H2O2, producing the weak oxidizing agent hypothiocyanite (OSCN−), which has bacteriostatic activity. TPO uses I− ions and H2O2 to generate iodine, and plays a central role in the biosynthesis of thyroid hormones T3 and T4. Myeloperoxidase (), for example, resides in the human nucleus and lysosome and acts as a defense response to oxidative stress, preventing apoptosis of the cell.
Structure
3D structures of MPO and PGHS have been reported. MPO is a homodimer: each monomer consists of a light (A or B) and a heavy (C or D) chain resulting from post-translational excision of 6 residues from the common precursor. Monomers are linked by a single inter-chain disulfide. Each monomer includes a bound calcium ion. PGHS exists as a symmetric homodimer, each monomer of which consists of 3 domains: an N-terminal epidermal growth factor (EGF) like module; a membrane-binding domain; and a large C-terminal catalytic domain containing the cyclooxygenase and the peroxidase active sites. The catalytic domain shows striking structural similarity to MPO. The image at the top of this page is an example of Myeloperoxidase 1dnu derived from X-ray diffraction with resolution 1.85 angstrom.
Active site
The cyclooxygenase active site, which catalyzes the formation of prostaglandin G2 (PGG2) from arachidonic acid, resides at the apex of a long hydrophobic channel, extending from the membrane-binding domain to the center of the molecule. The peroxidase active site, which catalyzes the reduction of PGG2 to PGH2, is located on the other side of the molecule, at the heme binding site. Both MPO and the catalytic domain of PGHS are mainly alpha-helical, 19 helices being identified as topologically and spatially equivalent; PGHS contains 5 additional N-terminal helices that have no equivalent in MPO. In both proteins, three Asn residues in each monomer are glycosylated.
Human proteins containing this domain
The following is a list of human proteins containing this domain:
DUOX1; DUOX2; EPX; LPO; MPO; PTGS1; PTGS2; PXDNL; TPO
References
External links
Protein domains
Protein families
Integral monotopic proteins
EC 1.11.1 | Animal heme-dependent peroxidases | [
"Biology"
] | 813 | [
"Protein families",
"Protein domains",
"Protein classification"
] |
14,711,012 | https://en.wikipedia.org/wiki/Endonuclease/Exonuclease/phosphatase%20family | Endonuclease/Exonuclease/phosphatase family is a structural domain found in the large family of proteins including magnesium dependent endonucleases and many phosphatases involved in intracellular signaling.
Examples
AP endonuclease proteins ,
DNase I proteins ,
Synaptojanin, an inositol-1,4,5-trisphosphate phosphatase
Sphingomyelinase
Nocturnin, an NADPH 2' phosphatase
Subfamilies
Inositol polyphosphate related phosphatase
Human proteins containing this domain
2'-PDE; 2-PDE; ANGEL1; ANGEL2; APEX1; APEX2; CCRN4L; CNOT6;
CNOT6L; DNASE1; DNASE1L1; DNASE1L2; DNASE1L3; INPP5A; INPP5B; INPP5D;
INPP5E; INPPL1; KIAA1706; OCRL; PIB5PA; SKIP; SMPD2; SMPD3;
SYNJ1; SYNJ2; TTRAP; Nocturnin;
Notes
References
Protein domains
Peripheral membrane proteins
EC 3.1.3 | Endonuclease/Exonuclease/phosphatase family | [
"Biology"
] | 270 | [
"Protein domains",
"Protein classification"
] |
14,712,087 | https://en.wikipedia.org/wiki/Major%20facilitator%20superfamily | The major facilitator superfamily (MFS) is a superfamily of membrane transport proteins that facilitate movement of small solutes across cell membranes in response to chemiosmotic gradients.
Function
The major facilitator superfamily (MFS) are membrane proteins which are expressed ubiquitously in all kingdoms of life for the import or export of target substrates. The MFS family was originally believed to function primarily in the uptake of sugars but subsequent studies revealed that drugs, metabolites, oligosaccharides, amino acids and oxyanions were all transported by MFS family members. These proteins energetically drive transport utilizing the electrochemical gradient of the target substrate (uniporter), or act as a cotransporter where transport is coupled to the movement of a second substrate.
Fold
The basic fold of the MFS transporter is built around 12, or in some cases, 14 transmembrane helices (TMH), with two 6- (or 7- ) helix bundles formed by the N and C terminal homologous domains of the transporter which are connected by an extended cytoplasmic loop. The two halves of the protein pack against each other in a clam-shell fashion, sealing via interactions at the ends of the transmembrane helices and extracellular loops. This forms a large aqueous cavity at the center of the membrane, which is alternatively open to the cytoplasm or periplasm/extracellular space. Lining this aqueous cavity are the amino-acids which bind the substrates and define transporter specificity. Many MFS transporters are thought to be dimers through in vitro and in vivo methods, with some evidence to suggest a functional role for this oligomerization.
Mechanism
The alternating-access mechanism thought to underlie the transport of most MFS transport is classically described as the "rocker-switch" mechanism. In this model, the transporter opens to either the extracellular space or cytoplasm and simultaneously seals the opposing face of the transporter, preventing a continuous pathway across the membrane. For example, in the best studied MFS transporter, LacY, lactose and protons typically bind from the periplasm to specific sites within the aqueous cleft. This drives closure of the extracellular face, and opening of the cytoplasmic side, allowing substrate into the cell. Upon substrate release, the transporter recycles to the periplasmic facing orientation.
Exporters and antiporters of the MFS family follow a similar reaction cycle, though exporters bind substrate in the cytoplasm and extrude it to the extracellular or periplasmic space, while antiporters bind substrate in both states to drive each conformational change. While most MFS structures suggest large, rigid body structural changes with substrate binding, the movements may be small in the cases of small substrates, such as the nitrate transporter NarK.
Transport
The generalized transport reactions catalyzed by MFS porters are:
Uniport: S (out) ⇌ S (in)
Symport: S (out) + [H+ or Na+] (out) ⇌ S (in) + [H+ or Na+] (in)
Antiport: S1 (out) + S2 (in) ⇌ S1 (in) + S2 (out) (S1 may be H+ or a solute)
Substrate specificity
Though initially identified as sugar transporters, a function conserved from prokaryotes to mammals, the MFS family is notable for the great diversity of substrates transported by the superfamily. These range from small oxyanions to large peptide fragments. Other MFS transporters are notable for a lack of selectivity, extruding broad classes of drugs and xenobiotics. This substrate specificity is largely determined by specific side chains which line the aqueous pocket at the center of the membrane. While one substrate of particular biological importance is often used to name the transporter or family, there may also be co-transported or leaked ions or molecules. These include water molecules or the coupling ions which energetically drive transport.
Structures
The crystal structures of a number of MFS transporters have been characterized. The first structures were of the glycerol 3-phosphate/phosphate exchanger GlpT and the lactose-proton symporter LacY, which served to elucidate the overall structure of the protein family and provided initial models for understanding the MFS transport mechanism. Since these initial structures other MFS structures have been solved which illustrate substrate specificity or states within the reaction cycle. While the initial MFS structures solved were of bacterial transporters, recently structures of the first eukaryotic structures have been published. These include a fungal phosphate transporter PiPT, plant nitrate transporter NRT1.1, and the human glucose transporter GLUT1.
Evolution
The origin of the basic MFS transporter fold is currently under heavy debate. All currently recognized MFS permeases have the two six-TMH domains within a single polypeptide chain, although in some MFS families an additional two TMHs are present. Evidence suggests that the MFS permeases arose by a tandem intragenic duplication event in the early prokaryotes. This event generated the 12 transmembrane helix topology from a (presumed) primordial 6-helix dimer. Moreover, the well-conserved MFS specific motif between TMS2 and TMS3 and the related but less well conserved motif between TMS8 and TMS9 prove to be a characteristic of virtually all of the more than 300 MFS proteins identified. However, the origin of the primordial 6-helix domain is under heavy debate. While some functional and structural evidence suggests that this domain arose out of a simpler 3-helix domain, bioinformatic or phylogenetic evidence supporting this hypothesis is lacking.
Medical significance
MFS family members are central to human physiology and play an important role in a number of diseases, through aberrant action, drug transport, or drug resistance. The OAT1 transporter transports a number of nucleoside analogs central to antiviral therapy. Resistance to antibiotics is frequently the result of action of MFS resistance genes. Mutations in MFS transporters have also been found to cause neurodegerative disease, vascular disorders of the brain, and glucose storage diseases.
Disease mutations
Disease associated mutations have been found in a number of human MFS transporters; those annotated in Uniprot are listed below.
Human MFS proteins
There are several MFS proteins in humans, where they are known as solute carriers (SLCs) and Atypical SLCs. There are today 52 SLC families, of which 16 families include MFS proteins; SLC2, 15 16, 17, 18, 19, SLCO (SLC21), 22, 29, 33, 37, 40, 43, 45, 46 and 49. Atypical SLCs are MFS proteins, sharing sequence similarities and evolutionary origin with SLCs, but they are not named according to the SLC root system, which originates from the hugo gene nomenclature system (HGNC). All atypical SLCs are listed in detail in, but they are: MFSD1, MFSD2A, MFSD2B, MFSD3, MFSD4A, MFSD4B, MFSD5, MFSD6, MFSD6L, MFSD8, MFSD9, MFSD10, MFSD11, MFSD12, MFSD13A, MFSD14A, MFSD14B, UNC93A, UNC93B1, SV2A, SV2B, SV2C, SVOP, SVOPL, SPNS1, SPNS2, SPNS3 and CLN3. As there is high sequence identity and phylogenetic resemblance between the atypical SLCs of MFS type, they can be divided into 15 AMTFs (Atypical MFS Transporter Families), suggesting there are at least 64 different families including SLC proteins of MFS type.
References
Protein domains
Transmembrane proteins
Articles containing video clips
Protein superfamilies
Transport proteins | Major facilitator superfamily | [
"Biology"
] | 1,735 | [
"Protein superfamilies",
"Protein domains",
"Protein classification"
] |
14,712,290 | https://en.wikipedia.org/wiki/Delta%20endotoxins | Delta endotoxins (δ-endotoxins) are a family of pore-forming toxins produced by Bacillus thuringiensis species of bacteria. They are useful for their insecticidal action and are the primary toxin produced by the genetically modified (GM) Bt maize/corn and other GM crops. During spore formation the bacteria produce crystals of such proteins (hence the name Cry toxins) that are also known as parasporal bodies, next to the endospores; as a result some members are known as a parasporin. The Cyt (cytolytic) toxin group is another group of delta-endotoxins formed in the cytoplasm. VIP toxins (vegetative insecticidal proteins) are formed at other stages of the life cycle.
Mechanism of action
When an insect ingests these proteins, they are activated by proteolytic cleavage. The N-terminus is cleaved in all of the proteins and a C-terminal extension is cleaved in some members. Once activated, the endotoxin binds to the gut epithelium and causes cell lysis by the formation of cation-selective channels, which leads to death.
For many years there was no clarity as to the relationship between aminopeptidase N and Bt toxins. Although AP-N does bind Cry proteins in vitro (reviewed by Soberón et al. 2009 and Pigott & Ellar 2007), no cases of resistance or even reduced in vitro binding due to AP-N structure alteration were known through 2002, and there was some doubt that the resistance mechanism was so straight forward. Indeed, Luo et al. 1997, Mohammed et al. 1996, and Zhu et al. 2000 positively found this to not occur in Lepidoptera examples. Subsequently, however Herrero et al. 2005 showed correlation between nonexpression and Bt resistance, and actual resistance was found in Helicoverpa armigera by Zhang et al. 2009, in Ostrinia nubilalis by Khajuria et al. 2011, and in Trichoplusia ni by Baxter et al. 2011 and Tiewsiri & Wang 2011 (also all Lepidoptera). There continues to be confirmation that AP-Ns do not by themselves affect resistance in some cases, possibly due to sequential binding by the toxin being required to produce its effect. In this sequence each binding step is theoretically not indispensable, but if it occurs does contribute to the final pore formation result.
Structure
The activated region of the delta toxin is composed of three distinct structural domains: an N-terminal helical bundle domain () involved in membrane insertion and pore formation; a beta-sheet central domain involved in receptor binding; and a C-terminal beta-sandwich domain () that interacts with the N-terminal domain to form a channel.
Types
B. thuringiensis encodes many proteins of the delta endotoxin family (), with some strains encoding multiple types simultaneously. A gene mostly found on plasmids, delta-entotoxins sometimes show up in genomes of other species, albeit at a lower proportion than those found in B. thuringiensis. The gene names looks like Cry3Bb, which in this case indicates a Cry toxin of superfamily 3 family B subfamily b.
Cry proteins that are interesting to cancer research are listed under a parasporin (PS) nomenclature in addition to the Cry nomenclature. They do not kill insects, but instead kill leukemia cells. The Cyt toxins tend to form their own group distinct from Cry toxins. Not all Cry crystal-form toxins directly share a common root. Examples of non-three-domain toxins that nevertheless have a Cry name include Cry34/35Ab1 and related beta-sandwich binary (Bin-like) toxins, Cry6Aa, and many beta-sandwich parasporins.
Specific delta-endotoxins that have been inserted with genetic engineering include Cry3Bb1 found in MON 863 and Cry1Ab found in MON 810, both of which are maize/corn cultivars. Cry3Bb1 is particularly useful because it kills Coleopteran insects such as the corn rootworm, an activity not seen in other Cry proteins. Other common toxins include Cry2Ab and Cry1F in cotton and maize/corn. In addition, Cry1Ac is effective as a vaccine adjuvant in humans.
Some insects populations have started to develop resistance towards delta endotoxin, with five resistant species found as of 2013. Plants with two kinds of delta endotoxins tend to make resistance happen slower, as the insects have to evolve to overcome both toxins at once. Planting non-Bt plants with the resistant plants will reduce the selection pressure for developing the toxin. Finally, two-toxin plants should not be planted with one-toxin plants, as one-toxin plants act as a stepping stone for adaption in this case.
References
Further reading
External links
Cry3Bb1 at the United States Environmental Protection Agency
Protein domains
Peripheral membrane proteins
Bacterial toxins
Crystals
Proteins | Delta endotoxins | [
"Chemistry",
"Materials_science",
"Biology"
] | 1,050 | [
"Biomolecules by chemical classification",
"Protein classification",
"Crystallography",
"Crystals",
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"Molecular biology",
"Proteins"
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1,039,124 | https://en.wikipedia.org/wiki/Stellar%20structure | Stellar structure models describe the internal structure of a star in detail and make predictions about the luminosity, the color and the future evolution of the star. Different classes and ages of stars have different internal structures, reflecting their elemental makeup and energy transport mechanisms.
Heat transport
For energy transport refer to Radiative transfer.
Different layers of the stars transport heat up and outwards in different ways, primarily convection and radiative transfer, but thermal conduction is important in white dwarfs.
Convection is the dominant mode of energy transport when the temperature gradient is steep enough so that a given parcel of gas within the star will continue to rise if it rises slightly via an adiabatic process. In this case, the rising parcel is buoyant and continues to rise if it is warmer than the surrounding gas; if the rising parcel is cooler than the surrounding gas, it will fall back to its original height. In regions with a low temperature gradient and a low enough opacity to allow energy transport via radiation, radiation is the dominant mode of energy transport.
The internal structure of a main sequence star depends upon the mass of the star.
In stars with masses of 0.3–1.5 solar masses (), including the Sun, hydrogen-to-helium fusion occurs primarily via proton–proton chains, which do not establish a steep temperature gradient. Thus, radiation dominates in the inner portion of solar mass stars. The outer portion of solar mass stars is cool enough that hydrogen is neutral and thus opaque to ultraviolet photons, so convection dominates. Therefore, solar mass stars have radiative cores with convective envelopes in the outer portion of the star.
In massive stars (greater than about 1.5 ), the core temperature is above about 1.8×107 K, so hydrogen-to-helium fusion occurs primarily via the CNO cycle. In the CNO cycle, the energy generation rate scales as the temperature to the 15th power, whereas the rate scales as the temperature to the 4th power in the proton-proton chains. Due to the strong temperature sensitivity of the CNO cycle, the temperature gradient in the inner portion of the star is steep enough to make the core convective. In the outer portion of the star, the temperature gradient is shallower but the temperature is high enough that the hydrogen is nearly fully ionized, so the star remains transparent to ultraviolet radiation. Thus, massive stars have a radiative envelope.
The lowest mass main sequence stars have no radiation zone; the dominant energy transport mechanism throughout the star is convection.
Equations of stellar structure
The simplest commonly used model of stellar structure is the spherically symmetric quasi-static model, which assumes that a star is in a steady state and that it is spherically symmetric. It contains four basic first-order differential equations: two represent how matter and pressure vary with radius; two represent how temperature and luminosity vary with radius.
In forming the stellar structure equations (exploiting the assumed spherical symmetry), one considers the matter density , temperature , total pressure (matter plus radiation) , luminosity , and energy generation rate per unit mass in a spherical shell of a thickness at a distance from the center of the star. The star is assumed to be in local thermodynamic equilibrium (LTE) so the temperature is identical for matter and photons. Although LTE does not strictly hold because the temperature a given shell "sees" below itself is always hotter than the temperature above, this approximation is normally excellent because the photon mean free path, , is much smaller than the length over which the temperature varies considerably, i.e. .
First is a statement of hydrostatic equilibrium: the outward force due to the pressure gradient within the star is exactly balanced by the inward force due to gravity. This is sometimes referred to as stellar equilibrium.
,
where is the cumulative mass inside the shell at and G is the gravitational constant. The cumulative mass increases with radius according to the mass continuity equation:
Integrating the mass continuity equation from the star center () to the radius of the star () yields the total mass of the star.
Considering the energy leaving the spherical shell yields the energy equation:
,
where is the luminosity produced in the form of neutrinos (which usually escape the star without interacting with ordinary matter) per unit mass. Outside the core of the star, where nuclear reactions occur, no energy is generated, so the luminosity is constant.
The energy transport equation takes differing forms depending upon the mode of energy transport. For conductive energy transport (appropriate for a white dwarf), the energy equation is
where k is the thermal conductivity.
In the case of radiative energy transport, appropriate for the inner portion of a solar mass main sequence star and the outer envelope of a massive main sequence star,
where is the opacity of the matter, is the Stefan–Boltzmann constant, and the Boltzmann constant is set to one.
The case of convective energy transport does not have a known rigorous mathematical formulation, and involves turbulence in the gas. Convective energy transport is usually modeled using mixing length theory. This treats the gas in the star as containing discrete elements which roughly retain the temperature, density, and pressure of their surroundings but move through the star as far as a characteristic length, called the mixing length. For a monatomic ideal gas, when the convection is adiabatic, meaning that the convective gas bubbles don't exchange heat with their surroundings, mixing length theory yields
where is the adiabatic index, the ratio of specific heats in the gas. (For a fully ionized ideal gas, .) When the convection is not adiabatic, the true temperature gradient is not given by this equation. For example, in the Sun the convection at the base of the convection zone, near the core, is adiabatic but that near the surface is not. The mixing length theory contains two free parameters which must be set to make the model fit observations, so it is a phenomenological theory rather than a rigorous mathematical formulation.
Also required are the equations of state, relating the pressure, opacity and energy generation rate to other local variables appropriate for the material, such as temperature, density, chemical composition, etc. Relevant equations of state for pressure may have to include the perfect gas law, radiation pressure, pressure due to degenerate electrons, etc. Opacity cannot be expressed exactly by a single formula. It is calculated for various compositions at specific densities and temperatures and presented in tabular form. Stellar structure codes (meaning computer programs calculating the model's variables) either interpolate in a density-temperature grid to obtain the opacity needed, or use a fitting function based on the tabulated values. A similar situation occurs for accurate calculations of the pressure equation of state. Finally, the nuclear energy generation rate is computed from nuclear physics experiments, using reaction networks to compute reaction rates for each individual reaction step and equilibrium abundances for each isotope in the gas.
Combined with a set of boundary conditions, a solution of these equations completely describes the behavior of the star. Typical boundary conditions set the values of the observable parameters appropriately at the surface () and center () of the star: , meaning the pressure at the surface of the star is zero; , there is no mass inside the center of the star, as required if the mass density remains finite; , the total mass of the star is the star's mass; and , the temperature at the surface is the effective temperature of the star.
Although nowadays stellar evolution models describe the main features of color–magnitude diagrams, important improvements have to be made in order to remove uncertainties which are linked to the limited knowledge of transport phenomena. The most difficult challenge remains the numerical treatment of turbulence. Some research teams are developing simplified modelling of turbulence in 3D calculations.
Rapid evolution
The above simplified model is not adequate without modification in situations when the composition changes are sufficiently rapid. The equation of hydrostatic equilibrium may need to be modified by adding a radial acceleration term if the radius of the star is changing very quickly, for example if the star is radially pulsating. Also, if the nuclear burning is not stable, or the star's core is rapidly collapsing, an entropy term must be added to the energy equation.
See also
Scale height
Standard solar model
References
Sources
External links
opacity code retrieved November 2009
The Yellow CESAM code, stellar evolution and structure Fortran source code
EZ to Evolve ZAMS Stars a FORTRAN 90 software derived from Eggleton's Stellar Evolution Code, a web-based interface can be found here .
Geneva Grids of Stellar Evolution Models (some of them including rotational induced mixing)
The BaSTI database of stellar evolution tracks
Stellar atmospheres: A contribution to the observational study of high temperature in the reversing layers of stars, (1925) by Cecilia Payne-Gaposchkin, Cambridge: The Observatory.
Structure
Stellar astronomy classification systems
Concepts in stellar astronomy | Stellar structure | [
"Physics",
"Astronomy"
] | 1,827 | [
"Stellar astronomy classification systems",
"Concepts in astrophysics",
"Astronomical classification systems",
"Concepts in stellar astronomy",
"Astronomical sub-disciplines",
"Stellar astronomy"
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1,039,260 | https://en.wikipedia.org/wiki/HOMFLY%20polynomial | In the mathematical field of knot theory, the HOMFLY polynomial or HOMFLYPT polynomial, sometimes called the generalized Jones polynomial, is a 2-variable knot polynomial, i.e. a knot invariant in the form of a polynomial of variables m and l.
A central question in the mathematical theory of knots is whether two knot diagrams represent the same knot. One tool used to answer such questions is a knot polynomial, which is computed from a diagram of the knot and can be shown to be an invariant of the knot, i.e. diagrams representing the same knot have the same polynomial. The converse may not be true. The HOMFLY polynomial is one such invariant and it generalizes two polynomials previously discovered, the Alexander polynomial and the Jones polynomial, both of which can be obtained by appropriate substitutions from HOMFLY. The HOMFLY polynomial is also a quantum invariant.
The name HOMFLY combines the initials of its co-discoverers: Jim Hoste, Adrian Ocneanu, Kenneth Millett, Peter J. Freyd, W. B. R. Lickorish, and David N. Yetter. The addition of PT recognizes independent work carried out by Józef H. Przytycki and Paweł Traczyk.
Definition
The polynomial is defined using skein relations:
where are links formed by crossing and smoothing changes on a local region of a link diagram, as indicated in the figure.
The HOMFLY polynomial of a link L that is a split union of two links and is given by
See the page on skein relation for an example of a computation using such relations.
Other HOMFLY skein relations
This polynomial can be obtained also using other skein relations:
Main properties
, where # denotes the knot sum; thus the HOMFLY polynomial of a composite knot is the product of the HOMFLY polynomials of its components.
, so the HOMFLY polynomial can often be used to distinguish between two knots of different chirality. However there exist chiral pairs of knots that have the same HOMFLY polynomial, e.g. knots 942 and 1071 together with their respective mirror images.
The Jones polynomial, V(t), and the Alexander polynomial, can be computed in terms of the HOMFLY polynomial (the version in and variables) as follows:
References
Further reading
Kauffman, L.H., "Formal knot theory", Princeton University Press, 1983.
Lickorish, W.B.R. "An Introduction to Knot Theory". Springer. .
External links
Knot theory
Polynomials | HOMFLY polynomial | [
"Mathematics"
] | 525 | [
"Polynomials",
"Algebra"
] |
1,039,609 | https://en.wikipedia.org/wiki/Antimatter%20weapon | An antimatter weapon is a theoretically possible device using antimatter as a power source, a propellant, or an explosive for a weapon. Antimatter weapons are currently too costly and unreliable to be viable in warfare, as producing antimatter is enormously expensive (estimated at US$6 billion for every 100 nanograms), the quantities of antimatter generated are very small, and current technology has great difficulty containing antimatter, which annihilates upon touching ordinary matter.
The paramount advantage of such a theoretical weapon is that antimatter and matter collisions result in the entire sum of their mass energy equivalent being released as energy, which is at least two orders of magnitude greater than the energy release of the most efficient fusion weapons (100% vs 0.4–1%). Annihilation requires and converts exactly equal masses of antimatter and matter by the collision which releases the entire mass-energy of both, which for 1 gram is ~9×1013 joules. Using the convention that 1 kiloton TNT equivalent = 4.184×1012 joules (or one trillion calories of energy), one half gram of antimatter reacting with one half gram of ordinary matter (one gram total) results in 21.5 kilotons-equivalent of energy (the same as the atomic bomb dropped on Nagasaki in 1945).
Cost
, the cost of producing one millionth of a gram of antimatter was estimated at . By way of comparison, the cost of the Manhattan Project (to produce the first atomic bomb) is estimated at US$23 billion in 2007 prices. As such, Hui Chen of Lawrence Livermore National Laboratory dismissed concerns about antimatter bombs in 2008 as "unrealistic".
Antimatter catalyzed weapons
Antimatter-catalyzed nuclear pulse propulsion proposes the use of antimatter as a "trigger" to initiate small nuclear explosions; the explosions provide thrust to a spacecraft. The same technology could theoretically be used to make very small and possibly "fission-free" (very low nuclear fallout) weapons (see pure fusion weapon).
In popular culture
An antimatter weapon is a part of the plot of the Dan Brown book Angels & Demons and its film adaptation, where it is used in a plot to blow up the Vatican City.
The Ground Zero expansion pack of the video game Quake II requires the protagonist to manufacture an Antimatter Bomb in the Munitions Plant to achieve the final objective.
In the Star Trek franchise, Federation starships are armed with photon torpedoes which contain antimatter warheads.
References
External links
Spotlight on "Angels and Demons" – A discussion at CERN's public website on the viability of the use of antimatter for energy and weaponry
"Air Force pursuing antimatter weapons: Program was touted publicly, then came official gag order"
Page discussing the possibility of using antimatter as a trigger for a thermonuclear explosion
Paper discussing the number of antiprotons required to ignite a thermonuclear weapon
Weapon
Proposed weapons
Science fiction weapons
Weapons of mass destruction | Antimatter weapon | [
"Physics"
] | 636 | [
"Antimatter",
"Matter"
] |
1,039,777 | https://en.wikipedia.org/wiki/Vis-viva%20equation | In astrodynamics, the vis-viva equation is one of the equations that model the motion of orbiting bodies. It is the direct result of the principle of conservation of mechanical energy which applies when the only force acting on an object is its own weight which is the gravitational force determined by the product of the mass of the object and the strength of the surrounding gravitational field.
Vis viva (Latin for "living force") is a term from the history of mechanics and this name is given to the orbital equation originally derived by Isaac Newton. It represents the principle that the difference between the total work of the accelerating forces of a system and that of the retarding forces is equal to one half the vis viva accumulated or lost in the system while the work is being done.
Equation
For any Keplerian orbit (elliptic, parabolic, hyperbolic, or radial), the vis-viva equation is as follows:
where:
is the relative speed of the two bodies
is the distance between the two bodies' centers of mass
is the length of the semi-major axis ( for ellipses, or for parabolas, and for hyperbolas)
is the gravitational constant
is the mass of the central body
The product of can also be expressed as the standard gravitational parameter using the Greek letter .
Practical applications
Given the total mass and the scalars and at a single point of the orbit, one can compute:
and at any other point in the orbit; and
the specific orbital energy , allowing an object orbiting a larger object to be classified as having not enough energy to remain in orbit, hence being "suborbital" (a ballistic missile, for example), having enough energy to be "orbital", but without the possibility to complete a full orbit anyway because it eventually collides with the other body, or having enough energy to come from and/or go to infinity (as a meteor, for example).
The formula for escape velocity can be obtained from the Vis-viva equation by taking the limit as approaches :
For a given orbital radius, the escape velocity will be times the orbital velocity.
Derivation for elliptic orbits (0 ≤ eccentricity < 1)
Specific total energy is constant throughout the orbit. Thus, using the subscripts and to denote apoapsis (apogee) and periapsis (perigee), respectively,
Rearranging,
Recalling that for an elliptical orbit (and hence also a circular orbit) the velocity and radius vectors are perpendicular at apoapsis and periapsis, conservation of angular momentum requires specific angular momentum , thus :
Isolating the kinetic energy at apoapsis and simplifying,
From the geometry of an ellipse, where a is the length of the semimajor axis. Thus,
Substituting this into our original expression for specific orbital energy,
Thus, and the vis-viva equation may be written
or
Therefore, the conserved angular momentum can be derived using and , where is semi-major axis and is semi-minor axis of the elliptical orbit, as follows:
and alternately,
Therefore, specific angular momentum , and
Total angular momentum
References
Orbits
Conservation laws
Equations of astronomy | Vis-viva equation | [
"Physics",
"Astronomy"
] | 635 | [
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"Physics theorems"
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1,039,945 | https://en.wikipedia.org/wiki/Royal%20Meteorological%20Society | The Royal Meteorological Society is a long-established institution that promotes academic and public engagement in weather and climate science. Fellows of the Society must possess relevant qualifications, but Members can be lay enthusiasts. Its Quarterly Journal is one of the world's leading sources of original research in the atmospheric sciences. The chief executive officer is Liz Bentley.
Constitution
The Royal Meteorological Society traces its origins back to 3 April 1850 when the British Meteorological Society was formed as "a society the objects of which should be the advancement and extension of meteorological science by determining the laws of climate and of meteorological phenomena in general". Along with nine others, including James Glaisher, John Drew, Edward Joseph Lowe, The Revd Joseph Bancroft Reade, and Samuel Charles Whitbread, Dr John Lee, an astronomer, of Hartwell House, near Aylesbury, Buckinghamshire founded in the library of his house the British Meteorological Society, which became the Royal Meteorological Society. It became The Meteorological Society in 1866, when it was incorporated by Royal Charter, and the Royal Meteorological Society in 1883, when Her Majesty Queen Victoria granted the privilege of adding 'Royal' to the title. Along with 74 others, the famous meteorologist Luke Howard joined the original 15 members of the Society at its first ordinary meeting on 7 May 1850. As of 2008 it has more than 3,000 members worldwide. The chief executive of the Society is Professor Liz Bentley. Paul Hardaker previously served as chief executive from 2006 to 2012.
Membership
There are four membership categories:
Honorary Fellow
Fellow (FRMetS)
Member
Corporate member
Awards
The society regularly awards a number of medal and prizes, of which the Symons Gold Medal (established in 1901) and the Mason Gold Medal (established in 2006) are pre-eminent. The two medals are awarded alternately.
Other awards include the Buchan Prize, the Hugh Robert Mill Award, the L F Richardson Prize, the Michael Hunt Award, the Fitzroy Prize, the Gordon Manley Weather Prize, the International Journal of Climatology Prize, the Society Outstanding Service Award and the Vaisala Award.
Journals
The society has a number of regular publications:
Atmospheric Science Letters: a monthly journal that provides a peer-reviewed publication route for new shorter contributions in the field of atmospheric and closely related sciences.
Weather: a monthly journal with many full colour illustrations and photos for specialists and general readers with an interest in meteorology. It uses a minimum of mathematics and technical language.
Quarterly Journal of the Royal Meteorological Society: one of the world's leading journals for meteorology, publishing original research in the atmospheric sciences. There are eight issues per year.
Meteorological Applications: this is a journal for applied meteorologists, forecasters and users of meteorological services and has been published since 1994. It is aimed at a general readership and authors are asked to take this into account when preparing papers.
International Journal of Climatology: has 15 issues a year and covers a broad spectrum of research in climatology.
WIREs Climate Change: a journal about climate change
Geoscience Data Journal: an online, open-access journal.
Climate Resilience and Sustainability: an interdisciplinary, open-access journal.
All publications are available online but a subscription is required for some. However certain "classic" papers are freely available on the Society's website.
Local centres and special interest groups
The society has several local centres across the UK.
There are also a number of special interest groups which organise meetings and other activities to facilitate exchange of information and views within specific areas of meteorology. These are informal groups of professionals interested in specific technical areas of the profession of meteorology. The groups are primarily a way of communicating at a specialist level.
Presidents
Source:
1850–1853: Samuel Charles Whitbread, first time
1853–1855: George Leach
1855–1857: John Lee
1857–1858: Robert Stephenson
1859–1860: Thomas Sopwith
1861–1862: Nathaniel Beardmore
1863–1864: Robert Dundas Thomson, died in office
1864: Samuel Charles Whitbread, second time
1865–1866: Charles Brooke
1867–1868: James Glaisher
1869–1870: Charles Vincent Walker
1871–1872: John William Tripe
1873–1875: Robert James Mann
1876–1877: Henry Storks Eaton
1878–1879: Charles Greaves
1880–1881: George James Symons, first time
1882–1883: Sir John Knox Laughton
1884–1885: Robert Henry Scott
1886–1887: William Ellis
1888–1889: William Marcet
1890–1891: Baldwin Latham
1892–1893: Charles Theodore Williams, first time
1894–1895: Richard Inwards
1896–1897: Edward Mawley
1898–1899: Francis Campbell Bayard
1900: George James Symons, second time; died in office
1900: Charles Theodore Williams, second time
1901–1902: William Henry Dines
1903–1904: Captain David W. Barker
1905–1906: Richard Bentley
1907–1908: Hugh Robert Mill
1910–1911: Henry Mellish
1911–1912: Henry Newton Dickson
1913–1914: Charles John Philip Cave, first time
1915–1917: Sir Henry George Lyons
1918–1919: Sir Napier Shaw
1920–1921: Reginald Hawthorn Hooker
1922–1923: Charles Chree
1924–1925: Charles John Philip Cave, second time
1926–1927: Sir Gilbert Walker
1928–1929: Richard Gregory
1930–1931: Rudolf Gustav Karl Lempfert
1932–1933: Sydney Chapman
1934–1935: Ernest Gold
1936–1937: Francis John Welsh Whipple
1938–1939: Sir Bernard A. Keen
1940–1941: Sir George Clarke Simpson
1942–1944: David Brunt
1945–1946: Gordon Manley
1947–1949: G. M. B. Dobson
1949–1951: Sir Robert Alexander Watson-Watt
1951–1953: Sir Charles Normand
1953–1955: Sir Graham Sutton
1955–1957: Reginald Sutcliffe
1957–1959: Percival Albert Sheppard
1959–1961: James Martin Stagg
1961–1963: Howard Latimer Penman
1963–1965: John Stanley Sawyer
1965–1967: G. D. Robinson
1967–1968: F. Kenneth Hare
1968–1970: John Mason
1970–1972: Frank Pasquill
1972–1974: Robert B. Pearce
1974–1976: Raymond Hide
1976–1978: John T. Houghton
1978–1980: John Monteith
1980–1982: Philip Goldsmith
1982–1984: Henry Charnock
1984–1986: Andrew Gilchrist
1986–1988: Richard S. Scorer
1988–1990: Keith Anthony Browning
1990–1992: Stephen Austen Thorpe
1992–1994: Paul James Mason
1994–1996: John E. Harries
1996–1998: David J. Carson
1998–2000: Sir Brian Hoskins
2000–2002: David Burridge
2002–2004: Howard Cattle
2004–2006: Chris Collier
2006–2008: Geraint Vaughan
2008–2010: Julia Slingo
2010–2012: Tim Palmer
2012–2014: Joanna Haigh
2014–2016: Jennie Campbell
2016–2018: Ellie Highwood
2018–2020: David Warrilow
2020–2022: David Griggs
Notable fellows
John Farrah (1849–1907).
See also
List of atmospheric dispersion models
UK Dispersion Modelling Bureau
Met Office
References
External links
The RMetS website
UK Atmospheric Dispersion Modelling Liaison Committee (ADMLC) web site
Meteorological societies
Meteorological
Scientific organisations based in the United Kingdom
Atmospheric dispersion modeling
Climatological research organizations
Climate of the United Kingdom
Geographic societies
Learned societies of the United Kingdom
Scientific organizations established in 1850
1850 establishments in the United Kingdom | Royal Meteorological Society | [
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1,040,853 | https://en.wikipedia.org/wiki/Radio%20window | The radio window is the region of the radio spectrum that penetrate the Earth's atmosphere. Typically, the lower limit of the radio window's range has a value of about 10 MHz (λ ≈ 30 m); the best upper limit achievable from optimal terrestrial observation sites is equal to approximately 1 THz (λ ≈ 0.3 mm).
It plays an important role in astronomy; up until the 1940s, astronomers could only use the visible and near infrared spectra for their measurements and observations. With the development of radio telescopes, the radio window became more and more utilizable, leading to the development of radio astronomy that provided astrophysicists with valuable observational data.
Factors affecting lower and upper limits
The lower and upper limits of the radio window's range of frequencies are not fixed; they depend on a variety of factors.
Absorption of mid-IR
The upper limit is affected by the vibrational transitions of atmospheric molecules such as oxygen (O2), carbon dioxide (CO2), and water (H2O), whose energies are comparable to the energies of mid-infrared photons: these molecules largely absorb the mid-infrared radiation that heads towards Earth.
Ionosphere
The radio window's lower frequency limit is greatly affected by the ionospheric refraction of the radio waves whose frequencies are approximately below 30 MHz (λ > 10 m); radio waves with frequencies below the limit of 10 MHz (λ > 30 m) are reflected back into space by the ionosphere. The lower limit is proportional to the density of the ionosphere's free electrons and coincides with the plasma frequency:
where is the plasma frequency in Hz and the electron density in electrons per cubic meter. Since it is highly dependent on sunlight, the value of changes significantly from daytime to nighttime usually being lower during the day, leading to a decrease of the radio window's lower limit and higher during the night, causing an increase of the radio window's lower frequency end. However, this also depends on the solar activity and the geographic position.
Troposphere
When performing observations, radio astronomers try to extend the upper limit of the radio window towards the 1 THz optimum, since the astronomical objects give spectral lines of greater intensity in the higher frequency range. Tropospheric water vapour greatly affects the upper limit since its resonant absorption frequency bands are 22.3 GHz (λ ≈ 1.32 cm), 183.3 GHz (λ ≈ 1.64 mm) and 323.8 GHz (λ ≈ 0.93 mm). The tropospheric oxygen's bands at 60 GHz (λ ≈ 5.00 mm) and 118.74 GHz (λ ≈ 2.52 mm) also affect the upper limit. To tackle the issue of water vapour, many observatories are built at high altitudes where the climate is more dry. However, few can be done to avoid the oxygen's interference with radio waves propagation.
Radio frequency interference
The width of the radio window is also affected by radio frequency interference which hinders the observations at certain wavelength ranges and undermines the quality of the observational data of radio astronomy.
See also
Infrared window
Optical window
Radio propagation
References
Electromagnetic spectrum | Radio window | [
"Physics"
] | 656 | [
"Spectrum (physical sciences)",
"Electromagnetic spectrum"
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1,040,953 | https://en.wikipedia.org/wiki/Recuperator | A recuperator is a special purpose counter-flow energy recovery heat exchanger positioned within the supply and exhaust air streams of an air handling system, or in the exhaust gases of an industrial process, in order to recover the waste heat. Generally, they are used to extract heat from the exhaust and use it to preheat air entering the combustion system. In this way they use waste energy to heat the air, offsetting some of the fuel, and thereby improve the energy efficiency of the system as a whole.
Description
In many types of processes, combustion is used to generate heat, and the recuperator serves to recuperate, or reclaim this heat, in order to reuse or recycle it. The term recuperator refers as well to liquid-liquid counterflow heat exchangers used for heat recovery in the chemical and refinery industries and in closed processes such as ammonia-water or LiBr-water absorption refrigeration cycle.
Recuperators are often used in association with the burner portion of a heat engine, to increase the overall efficiency. For example, in a gas turbine engine, air is compressed, mixed with fuel, which is then burned and used to drive a turbine. The recuperator transfers some of the waste heat in the exhaust to the compressed air, thus preheating it before entering the fuel burner stage. Since the gases have been pre-heated, less fuel is needed to heat the gases up to the turbine inlet temperature. By recovering some of the energy usually lost as waste heat, the recuperator can make a heat engine or gas turbine significantly more efficient.
Energy transfer process
Normally the heat transfer between airstreams provided by the device is termed as "sensible heat", which is the exchange of energy, or enthalpy, resulting in a change in temperature of the medium (air in this case), but with no change in moisture content. However, if moisture or relative humidity levels in the return air stream are high enough to allow condensation to take place in the device, then this will cause "latent heat" to be released and the heat transfer material will be covered with a film of water. Despite a corresponding absorption of latent heat, as some of the water film is evaporated in the opposite airstream, the water will reduce the thermal resistance of the boundary layer of the heat exchanger material and thus improve the heat transfer coefficient of the device, and hence increase efficiency. The energy exchange of such devices now comprises both sensible and latent heat transfer; in addition to a change in temperature, there is also a change in moisture content of the exhaust air stream.
However, the film of condensation will also slightly increase pressure drop through the device, and depending upon the spacing of the matrix material, this can increase resistance by up to 30%. If the unit is not laid to falls, and the condensate not allowed to drain properly, this will increase fan energy consumption and reduce the seasonal efficiency of the device.
Use in ventilation systems
In heating, ventilation and air-conditioning systems, HVAC, recuperators are commonly used to re-use waste heat from exhaust air normally expelled to atmosphere. Devices typically comprises a series of parallel plates of aluminium, plastic, stainless steel, or synthetic fiber, copper alternate pairs of which are enclosed on two sides to form twin sets of ducts at right angles to each other, and which contain the supply and extract air streams. In this manner heat from the exhaust air stream is transferred through the separating plates, and into the supply air stream. Manufacturers claim gross efficiencies of up to 95% depending upon the specification of the unit.
The characteristics of this device are attributable to the relationship between the physical size of the unit, in particular the air path distance, and the spacing of the plates. For an equal air pressure drop through the device, a small unit will have a narrower plate spacing and a lower air velocity than a larger unit, but both units may be just as efficient. Because of the cross-flow design of the unit, its physical size will dictate the air path length, and as this increases, heat transfer will increase but pressure drop will also increase, and so plate spacing is increased to reduce pressure drop, but this in turn will reduce heat transfer.
As a general rule a recuperator selected for a pressure drop of between will have a good efficiency, while having a small effect on fan power consumption, but will have in turn a higher seasonal efficiency than that for physically smaller, but higher pressure drop recuperator.
When heat recovery is not required, it is typical for the device to be bypassed by use of dampers arranged within the ventilation distribution system. Assuming the fans are fitted with inverter speed controls, set to maintain a constant pressure in the ventilation system, then the reduced pressure drop leads to a slowing of the fan motor and thus reducing power consumption, and in turn improves the seasonal efficiency of the system.
Use in metallurgical furnaces
Recuperators have also been used to recover heat from waste gasses to preheat combustion air and fuel for many years by metallic recuperators to reduce energy costs and carbon footprint of operation. Compared to alternatives such as regenerative furnaces, initial costs are lesser, there are no valves to be switching back and forth, there are no induced-draft fans and it does not require a web of gas ducts spread up all over the furnace.
Historically the recovery ratios of recuperators compared to regenerative burners were low. However, recent improvements to technology have allowed recuperators to recover 70-80% of the waste heat and pre-heated air up to is now possible.
Gas turbines
Recuperators can be used to increase the efficiency of gas turbines for power generation, provided the exhaust gas is hotter than the compressor outlet temperature. The exhaust heat from the turbine is used to pre-heat the air from the compressor before further heating in the combustor, reducing the fuel input required. The larger the temperature difference between turbine out and compressor out, the greater the benefit from the recuperator. Therefore, microturbines (<1 MW), which typically have low pressure ratios, have the most to gain from the use of a recuperator. In practice, a doubling of efficiency is possible through the use of a recuperator. The major practical challenge for a recuperator in microturbine applications is coping with the exhaust gas temperature, which can exceed .
Other types of gas-to-gas heat exchangers
Heat pipe
Run-around coil
Thermal wheel, or rotary heat exchanger (including enthalpy wheel and desiccant wheel)
Convection recuperator
Radiation recuperator
See also
Air handler
Energy recovery ventilation
Heat recovery ventilation
HVAC (heating, ventilation, and air conditioning)
Indoor air quality
Regenerative heat exchanger
Thermal comfort
References
External links
Energy conservation
Energy recovery
Engineering thermodynamics
Heat exchangers
Heat transfer
Heating, ventilation, and air conditioning
Heating
Industrial equipment
Low-energy building
Mechanical engineering
Sustainable building | Recuperator | [
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1,041,641 | https://en.wikipedia.org/wiki/Radiation%20hardening | Radiation hardening is the process of making electronic components and circuits resistant to damage or malfunction caused by high levels of ionizing radiation (particle radiation and high-energy electromagnetic radiation), especially for environments in outer space (especially beyond low Earth orbit), around nuclear reactors and particle accelerators, or during nuclear accidents or nuclear warfare.
Most semiconductor electronic components are susceptible to radiation damage, and radiation-hardened (rad-hard) components are based on their non-hardened equivalents, with some design and manufacturing variations that reduce the susceptibility to radiation damage. Due to the low demand and the extensive development and testing required to produce a radiation-tolerant design of a microelectronic chip, the technology of radiation-hardened chips tends to lag behind the most recent developments. They also typically cost more than their commercial counterparts.
Radiation-hardened products are typically tested to one or more resultant-effects tests, including total ionizing dose (TID), enhanced low dose rate effects (ELDRS), neutron and proton displacement damage, and single event effects (SEEs).
Problems caused by radiation
Environments with high levels of ionizing radiation create special design challenges. A single charged particle can knock thousands of electrons loose, causing electronic noise and signal spikes. In the case of digital circuits, this can cause results which are inaccurate or unintelligible. This is a particularly serious problem in the design of satellites, spacecraft, future quantum computers, military aircraft, nuclear power stations, and nuclear weapons. In order to ensure the proper operation of such systems, manufacturers of integrated circuits and sensors intended for the military or aerospace markets employ various methods of radiation hardening. The resulting systems are said to be rad(iation)-hardened, rad-hard, or (within context) hardened.
Major radiation damage sources
Typical sources of exposure of electronics to ionizing radiation are the Van Allen radiation belts for satellites, nuclear reactors in power plants for sensors and control circuits, particle accelerators for control electronics (particularly particle detector devices), residual radiation from isotopes in chip packaging materials, cosmic radiation for spacecraft and high-altitude aircraft, and nuclear explosions for potentially all military and civilian electronics.
Secondary particles result from interaction of other kinds of radiation with structures around the electronic devices.
Van Allen radiation belts contain electrons (up to about 10 MeV) and protons (up to 100s MeV) trapped in the geomagnetic field. The particle flux in the regions farther from the Earth can vary wildly depending on the actual conditions of the Sun and the magnetosphere. Due to their position they pose a concern for satellites.
Nuclear reactors produce gamma radiation and neutron radiation which can affect sensor and control circuits in nuclear power plants.
Particle accelerators produce high energy protons and electrons, and the secondary particles produced by their interactions produce significant radiation damage on sensitive control and particle detector components, of the order of magnitude of 10 MRad[Si]/year for systems such as the Large Hadron Collider.
Chip packaging materials were an insidious source of radiation that was found to be causing soft errors in new DRAM chips in the 1970s. Traces of radioactive elements in the packaging of the chips were producing alpha particles, which were then occasionally discharging some of the capacitors used to store the DRAM data bits. These effects have been reduced today by using purer packaging materials, and employing error-correcting codes to detect and often correct DRAM errors.
Cosmic rays come from all directions and consist of approximately 85% protons, 14% alpha particles, and 1% heavy ions, together with X-ray and gamma-ray radiation. Most effects are caused by particles with energies between 0.1 and 20 GeV. The atmosphere filters most of these, so they are primarily a concern for spacecraft and high-altitude aircraft, but can also affect ordinary computers on the surface.
Solar particle events come from the direction of the sun and consist of a large flux of high-energy (several GeV) protons and heavy ions, again accompanied by X-ray radiation.
Nuclear explosions produce a short and extremely intense surge through a wide spectrum of electromagnetic radiation, an electromagnetic pulse (EMP), neutron radiation, and a flux of both primary and secondary charged particles. In case of a nuclear war they pose a potential concern for all civilian and military electronics.
Radiation effects on electronics
Fundamental mechanisms
Two fundamental damage mechanisms take place:
Lattice displacement
Lattice displacement is caused by neutrons, protons, alpha particles, heavy ions, and very high energy gamma photons. They change the arrangement of the atoms in the crystal lattice, creating lasting damage, and increasing the number of recombination centers, depleting the minority carriers and worsening the analog properties of the affected semiconductor junctions. Counterintuitively, higher doses over a short time cause partial annealing ("healing") of the damaged lattice, leading to a lower degree of damage than with the same doses delivered in low intensity over a long time (LDR or Low Dose Rate). This type of problem is particularly significant in bipolar transistors, which are dependent on minority carriers in their base regions; increased losses caused by recombination cause loss of the transistor gain (see neutron effects). Components certified as ELDRS (Enhanced Low Dose Rate Sensitive)-free do not show damage with fluxes below 0.01 rad(Si)/s = 36 rad(Si)/h.
Ionization effects
Ionization effects are caused by charged particles, including ones with energy too low to cause lattice effects. The ionization effects are usually transient, creating glitches and soft errors, but can lead to destruction of the device if they trigger other damage mechanisms (e.g., a latchup). Photocurrent caused by ultraviolet and X-ray radiation may belong to this category as well. Gradual accumulation of holes in the oxide layer in MOSFET transistors leads to worsening of their performance, up to device failure when the dose is high enough (see total ionizing dose effects).
The effects can vary wildly depending on all the parameters – type of radiation, total dose and radiation flux, combination of types of radiation, and even the kind of device load (operating frequency, operating voltage, actual state of the transistor during the instant it is struck by the particle) – which makes thorough testing difficult, time-consuming, and requiring many test samples.
Resultant effects
The "end-user" effects can be characterized in several groups:
A neutron interacting with a semiconductor lattice will displace the atoms in the lattice. This leads to an increase in the count of recombination centers and deep-level defects, reducing the lifetime of minority carriers, thus affecting bipolar devices more than CMOS ones. Bipolar devices on silicon tend to show changes in electrical parameters at levels of 1010 to 1011 neutrons/cm2, while CMOS devices aren't affected until 1015 neutrons/cm2. The sensitivity of devices may increase together with increasing level of integration and decreasing size of individual structures. There is also a risk of induced radioactivity caused by neutron activation, which is a major source of noise in high energy astrophysics instruments. Induced radiation, together with residual radiation from impurities in component materials, can cause all sorts of single-event problems during the device's lifetime. GaAs LEDs, common in optocouplers, are very sensitive to neutrons. The lattice damage influences the frequency of crystal oscillators. Kinetic energy effects (namely lattice displacement) of charged particles belong here too.
Total ionizing dose effects
Total ionizing dose effects represent the cumulative damage of the semiconductor lattice (lattice displacement damage) caused by exposure to ionizing radiation over time. It is measured in rads and causes slow gradual degradation of the device's performance. A total dose greater than 5000 rads delivered to silicon-based devices in a timespan on the order of seconds to minutes will cause long-term degradation. In CMOS devices, the radiation creates electron–hole pairs in the gate insulation layers, which cause photocurrents during their recombination, and the holes trapped in the lattice defects in the insulator create a persistent gate biasing and influence the transistors' threshold voltage, making the N-type MOSFET transistors easier and the P-type ones more difficult to switch on. The accumulated charge can be high enough to keep the transistors permanently open (or closed), leading to device failure. Some self-healing takes place over time, but this effect is not too significant. This effect is the same as hot carrier degradation in high-integration high-speed electronics. Crystal oscillators are somewhat sensitive to radiation doses, which alter their frequency. The sensitivity can be greatly reduced by using swept quartz. Natural quartz crystals are especially sensitive. Radiation performance curves for TID testing may be generated for all resultant effects testing procedures. These curves show performance trends throughout the TID test process and are included in the radiation test report.
Transient dose effects
Transient dose effects result from a brief high-intensity pulse of radiation, typically occurring during a nuclear explosion. The high radiation flux creates photocurrents in the entire body of the semiconductor, causing transistors to randomly open, changing logical states of flip-flops and memory cells. Permanent damage may occur if the duration of the pulse is too long, or if the pulse causes junction damage or a latchup. Latchups are commonly caused by the X-rays and gamma radiation flash of a nuclear explosion. Crystal oscillators may stop oscillating for the duration of the flash due to prompt photoconductivity induced in quartz.
Systems-generated EMP effects
SGEMP effects are caused by the radiation flash traveling through the equipment and causing local ionization and electric currents in the material of the chips, circuit boards, electrical cables and cases.
Digital damage: SEE
Single-event effects (SEE) have been studied extensively since the 1970s. When a high-energy particle travels through a semiconductor, it leaves an ionized track behind. This ionization may cause a highly localized effect similar to the transient dose one - a benign glitch in output, a less benign bit flip in memory or a register or, especially in high-power transistors, a destructive latchup and burnout. Single event effects have importance for electronics in satellites, aircraft, and other civilian and military aerospace applications. Sometimes, in circuits not involving latches, it is helpful to introduce RC time constant circuits that slow down the circuit's reaction time beyond the duration of an SEE.
Single-event transient
An SET happens when the charge collected from an ionization event discharges in the form of a spurious signal traveling through the circuit. This is de facto the effect of an electrostatic discharge. it is considered a soft error, and is reversible.
Single-event upset
Single-event upsets (SEU) or transient radiation effects in electronics are state changes of memory or register bits caused by a single ion interacting with the chip. They do not cause lasting damage to the device, but may cause lasting problems to a system which cannot recover from such an error. It is otherwise a reversible soft error. In very sensitive devices, a single ion can cause a multiple-bit upset (MBU) in several adjacent memory cells. SEUs can become single-event functional interrupts (SEFI) when they upset control circuits, such as state machines, placing the device into an undefined state, a test mode, or a halt, which would then need a reset or a power cycle to recover.
Single-event latchup
An SEL can occur in any chip with a parasitic PNPN structure. A heavy ion or a high-energy proton passing through one of the two inner-transistor junctions can turn on the thyristor-like structure, which then stays "shorted" (an effect known as latch-up) until the device is power-cycled. As the effect can happen between the power source and substrate, destructively high current can be involved and the part may fail. This is a hard error, and is irreversible. Bulk CMOS devices are most susceptible.
Single-event snapback
A single-event snapback is similar to an SEL but not requiring the PNPN structure, and can be induced in N-channel MOS transistors switching large currents, when an ion hits near the drain junction and causes avalanche multiplication of the charge carriers. The transistor then opens and stays opened, a hard error which is irreversible.
Single-event induced burnout
An SEB may occur in power MOSFETs when the substrate right under the source region gets forward-biased and the drain-source voltage is higher than the breakdown voltage of the parasitic structures. The resulting high current and local overheating then may destroy the device. This is a hard error, and is irreversible.
Single-event gate rupture
SEGR are observed in power MOSFETs when a heavy ion hits the gate region while a high voltage is applied to the gate. A local breakdown then happens in the insulating layer of silicon dioxide, causing local overheating and destruction (looking like a microscopic explosion) of the gate region. It can occur even in EEPROM cells during write or erase, when the cells are subjected to a comparatively high voltage. This is a hard error, and is irreversible.
SEE testing
While proton beams are widely used for SEE testing due to availability, at lower energies proton irradiation can often underestimate SEE susceptibility. Furthermore, proton beams expose devices to risk of total ionizing dose (TID) failure which can cloud proton testing results or result in premature device failure. White neutron beams—ostensibly the most representative SEE test method—are usually derived from solid target-based sources, resulting in flux non-uniformity and small beam areas. White neutron beams also have some measure of uncertainty in their energy spectrum, often with high thermal neutron content.
The disadvantages of both proton and spallation neutron sources can be avoided by using mono-energetic 14 MeV neutrons for SEE testing. A potential concern is that mono-energetic neutron-induced single event effects will not accurately represent the real-world effects of broad-spectrum atmospheric neutrons. However, recent studies have indicated that, to the contrary, mono-energetic neutrons—particularly 14 MeV neutrons—can be used to quite accurately understand SEE cross-sections in modern microelectronics.
Radiation-hardening techniques
Physical
Hardened chips are often manufactured on insulating substrates instead of the usual semiconductor wafers. Silicon on insulator (SOI) and silicon on sapphire (SOS) are commonly used. While normal commercial-grade chips can withstand between 50 and 100 gray (5 and 10 krad), space-grade SOI and SOS chips can survive doses between 1000 and 3000 gray (100 and 300 krad). At one time many 4000 series chips were available in radiation-hardened versions (RadHard). While SOI eliminates latchup events, TID and SEE hardness are not guaranteed to be improved.
Choosing a substrate with wide band gap gives it higher tolerance to deep-level defects; e.g. silicon carbide or gallium nitride.
Use of a special process node provides increased radiation resistance. Due to the high development costs of new radiation hardened processes, the smallest "true" rad-hard (RHBP, Rad-Hard By Process) process is 150 nm as of 2016, however, rad-hard 65 nm FPGAs were available that used some of the techniques used in "true" rad-hard processes (RHBD, Rad-Hard By Design). As of 2019 110 nm rad-hard processes are available.
Bipolar integrated circuits generally have higher radiation tolerance than CMOS circuits. The low-power Schottky (LS) 5400 series can withstand 1000 krad, and many ECL devices can withstand 10,000 krad. Using edgeless CMOS transistors, which have an unconventional physical construction, together with an unconventional physical layout, can also be effective.
Magnetoresistive RAM, or MRAM, is considered a likely candidate to provide radiation hardened, rewritable, non-volatile conductor memory. Physical principles and early tests suggest that MRAM is not susceptible to ionization-induced data loss.
Capacitor-based DRAM is often replaced by more rugged (but larger, and more expensive) SRAM. SRAM cells have more transistors per cell than usual (which is 4T or 6T), which makes the cells more tolerant to SEUs at the cost of higher power consumption and size.
Shielding
Shielding the package against radioactivity is straightforward to reduce exposure of the bare device.
To protect against neutron radiation and the neutron activation of materials, it is possible to shield the chips themselves by use of depleted boron (consisting only of isotope boron-11) in the borophosphosilicate glass passivation layer protecting the chips, as naturally prevalent boron-10 readily captures neutrons and undergoes alpha decay (see soft error).
Logical
Error correcting code memory (ECC memory) uses redundant bits to check for and possibly correct corrupted data. Since radiation's effects damage the memory content even when the system is not accessing the RAM, a "scrubber" circuit must continuously sweep the RAM; reading out the data, checking the redundant bits for data errors, then writing back any corrections to the RAM.
Redundant elements can be used at the system level. Three separate microprocessor boards may independently compute an answer to a calculation and compare their answers. Any system that produces a minority result will recalculate. Logic may be added such that if repeated errors occur from the same system, that board is shut down.
Redundant elements may be used at the circuit level. A single bit may be replaced with three bits and separate "voting logic" for each bit to continuously determine its result (triple modular redundancy). This increases area of a chip design by a factor of 5, so must be reserved for smaller designs. But it has the secondary advantage of also being "fail-safe" in real time. In the event of a single-bit failure (which may be unrelated to radiation), the voting logic will continue to produce the correct result without resorting to a watchdog timer. System level voting between three separate processor systems will generally need to use some circuit-level voting logic to perform the votes between the three processor systems.
Hardened latches may be used.
A watchdog timer will perform a hard reset of a system unless some sequence is performed that generally indicates the system is alive, such as a write operation from an onboard processor. During normal operation, software schedules a write to the watchdog timer at regular intervals to prevent the timer from running out. If radiation causes the processor to operate incorrectly, it is unlikely the software will work correctly enough to clear the watchdog timer. The watchdog eventually times out and forces a hard reset to the system. This is considered a last resort to other methods of radiation hardening.
Military and space industry applications
Radiation-hardened and radiation tolerant components are often used in military and aerospace applications, including point-of-load (POL) applications, satellite system power supplies, step down switching regulators, microprocessors, FPGAs, FPGA power sources, and high efficiency, low voltage subsystem power supplies.
However, not all military-grade components are radiation hardened. For example, the US MIL-STD-883 features many radiation-related tests, but has no specification for single event latchup frequency. The Fobos-Grunt space probe may have failed due to a similar assumption.
The market size for radiation hardened electronics used in space applications was estimated to be $2.35 billion in 2021. A new study has estimated that this will reach approximately $4.76 billion by the year 2032.
Nuclear hardness for telecommunication
In telecommunication, the term nuclear hardness has the following meanings:
1) an expression of the extent to which the performance of a system, facility, or device is expected to degrade in a given nuclear environment, 2) the physical attributes of a system or electronic component that will allow survival in an environment that includes nuclear radiation and electromagnetic pulses (EMP).
Notes
Nuclear hardness may be expressed in terms of either susceptibility or vulnerability.
The extent of expected performance degradation (e.g., outage time, data lost, and equipment damage) must be defined or specified. The environment (e.g., radiation levels, overpressure, peak velocities, energy absorbed, and electrical stress) must be defined or specified.
The physical attributes of a system or component that will allow a defined degree of survivability in a given environment created by a nuclear weapon.
Nuclear hardness is determined for specified or actual quantified environmental conditions and physical parameters, such as peak radiation levels, overpressure, velocities, energy absorbed, and electrical stress. It is achieved through design specifications and it is verified by test and analysis techniques.
Examples of rad-hard computers
The System/4 Pi, made by IBM and used on board the Space Shuttle (AP-101 variant), is based on the System/360 architecture.
The RCA1802 8-bit CPU, introduced in 1976, was the first serially-produced radiation-hardened microprocessor.
PIC 1886VE, Russian 50 MHz microcontroller designed by Milandr and manufactured by Sitronics-Mikron on 180 nm bulk-silicon technology.
m68k based:
The Coldfire M5208 used by General Dynamics is a low power (1.5 W) radiation hardened alternative.
MIL-STD-1750A based:
The RH1750 manufactured by GEC-Plessey.
The Proton 100k SBC by Space Micro Inc., introduced in 2003, uses an updated voting scheme called TTMR which mitigates single event upset (SEU) in a single processor. The processor is Equator BSP-15.
The Proton200k SBC by Space Micro Inc, introduced in 2004, mitigates SEU with its patented time triple modular redundancy (TTMR) technology, and single event function interrupts (SEFI) with H-Core technology. The processor is the high speed Texas Instruments 320C6Xx series digital signal processor. The Proton200k operates at 4000 MIPS while mitigating SEU.
MIPS based:
The RH32 is produced by Honeywell Aerospace.
The Mongoose-V used by NASA is a 32-bit microprocessor for spacecraft onboard computer applications (i. e. New Horizons).
The KOMDIV-32 is a 32-bit microprocessor, compatible with MIPS R3000, developed by NIISI, manufactured by Kurchatov Institute, Russia.
PowerPC / POWER based:
The RAD6000 single-board computer (SBC), produced by BAE Systems, includes a rad-hard POWER1 CPU.
The RHPPC is produced by Honeywell Aerospace. Based on hardened PowerPC 603e.
The SP0 and SP0-S are produced by Aitech Defense Systems is a 3U cPCI SBC which utilizes the SOI PowerQUICC-III MPC8548E, PowerPC e500 based, capable of processing speeds ranging from 833 MHz to 1.18 GHz.
The RAD750 SBC, also produced by BAE Systems, and based on the PowerPC 750 processor, is the successor to the RAD6000.
The SCS750 built by Maxwell Technologies, which votes three PowerPC 750 cores against each other to mitigate radiation effects. Seven of those are used by the Gaia spacecraft.
The Boeing Company, through its Satellite Development Center, produces a radiation hardened space computer variant based on the PowerPC 750.
The BRE440 by Moog Inc. IBM PPC440 core based system-on-a-chip, 266 MIPS, PCI, 2x Ethernet, 2x UARTS, DMA controller, L1/L2 cache
The RAD5500 processor, is the successor to the RAD750 based on the PowerPC e5500.
SPARC based:
The ERC32 and LEON 2, 3, 4 and 5 are radiation hardened processors designed by Gaisler Research and the European Space Agency. They are described in synthesizable VHDL available under the GNU Lesser General Public License and GNU General Public License respectively.
The Gen 6 single-board computer (SBC), produced by Cobham Semiconductor Solutions (formerly Aeroflex Microelectronics Solutions), enabled for the LEON microprocessor.
ARM based:
The Vorago VA10820, a 32-bit ARMv6-M Cortex-M0.
NASA and the United States Air Force are developing HPSC, a Cortex-A53 based processor for future spacecraft use
ESA DAHLIA, a Cortex-R52 based processor
RISC-V based:
Cobham Gaisler NOEL-V 64-bit.
NASA Jet Propulsion Laboratory has selected Microchip Technology to develop a new HPSC processor, based on SiFive Intelligence X280
See also
Communications survivability
EMC-aware programming
Institute for Space and Defense Electronics, Vanderbilt University
Mars Reconnaissance Orbiter
MESSENGER Mercury probe
Mars rovers
Tempest (codename)
Juno Radiation Vault
References
Books and Reports
External links
Federal Standard 1037C (link )
(I)ntegrated Approach with COTS Creates Rad-Tolerant (SBC) for Space – By Chad Thibodeau, Maxwell Technologies; COTS Journal, Dec 2003
Sandia Labs to develop (...) radiation-hardened Pentium (...) for space and defense needs – Sandia press release, 8 Dec 1998(also includes a general "backgrounder" section on Sandia's manufacturing processes for radiation-hardening of microelectronics)
Radiation effects on quartz crystals
Vanderbilt University Institute for Space and Defense Electronics
Military communications
Integrated circuits
Avionics computers
Electronics manufacturing
Spaceflight
Radiation effects
Semiconductor device defects | Radiation hardening | [
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1,041,812 | https://en.wikipedia.org/wiki/Superpotential | In theoretical physics, the superpotential is a function in supersymmetric quantum mechanics. Given a superpotential, two "partner potentials" are derived that can each serve as a potential in the Schrödinger equation. The partner potentials have the same spectrum, apart from a possible eigenvalue of zero, meaning that the physical systems represented by the two potentials have the same characteristic energies, apart from a possible zero-energy ground state.
One-dimensional example
Consider a one-dimensional, non-relativistic particle with a two state internal degree of freedom called "spin". (This is not quite the usual notion of spin encountered in nonrelativistic quantum mechanics, because "real" spin applies only to particles in three-dimensional space.) Let b and its Hermitian adjoint b† signify operators which transform a "spin up" particle into a "spin down" particle and vice versa, respectively. Furthermore, take b and b† to be normalized such that the anticommutator {b,b†} equals 1, and take that b2 equals 0. Let p represent the momentum of the particle and x represent its position with [x,p]=i, where we use natural units so that . Let W (the superpotential) represent an arbitrary differentiable function of x and define the supersymmetric operators Q1 and Q2 as
The operators Q1 and Q2 are self-adjoint. Let the Hamiltonian be
where W''' signifies the derivative of W. Also note that {Q1,Q2}=0. Under these circumstances, the above system is a toy model of N=2 supersymmetry. The spin down and spin up states are often referred to as the "bosonic" and "fermionic" states, respectively, in an analogy to quantum field theory. With these definitions, Q1 and Q2 map "bosonic" states into "fermionic" states and vice versa. Restricting to the bosonic or fermionic sectors gives two partner potentials determined by
In four spacetime dimensions
In supersymmetric quantum field theories with four spacetime dimensions, which might have some connection to nature, it turns out that scalar fields arise as the lowest component of a chiral superfield, which tends to automatically be complex valued. We may identify the complex conjugate of a chiral superfield as an anti-chiral superfield. There are two possible ways to obtain an action from a set of superfields:
Integrate a superfield on the whole superspace spanned by and ,
or
Integrate a chiral superfield on the chiral half of a superspace, spanned by and , not on .
The second option tells us that an arbitrary holomorphic function of a set of chiral superfields can show up as a term in a Lagrangian which is invariant under supersymmetry. In this context, holomorphic means that the function can only depend on the chiral superfields, not their complex conjugates. We may call such a function W, the superpotential. The fact that W is holomorphic in the chiral superfields helps explain why supersymmetric theories are relatively tractable, as it allows one to use powerful mathematical tools from complex analysis. Indeed, it is known that W receives no perturbative corrections, a result referred to as the perturbative non-renormalization theorem. Note that non-perturbative processes may correct this, for example through contributions to the beta functions due to instantons.
See also
Komar superpotential
References
Stephen P. Martin, A Supersymmetry Primer''. .
B. Mielnik and O. Rosas-Ortiz, "Factorization: Little or great algorithm?", J. Phys. A: Math. Gen. 37: 10007-10035, 2004
Supersymmetry
Supersymmetric quantum field theory
Potentials | Superpotential | [
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1,042,053 | https://en.wikipedia.org/wiki/Neutron%20scattering | Neutron scattering, the irregular dispersal of free neutrons by matter, can refer to either the naturally occurring physical process itself or to the man-made experimental techniques that use the natural process for investigating materials. The natural/physical phenomenon is of elemental importance in nuclear engineering and the nuclear sciences. Regarding the experimental technique, understanding and manipulating neutron scattering is fundamental to the applications used in crystallography, physics, physical chemistry, biophysics, and materials research.
Neutron scattering is practiced at research reactors and spallation neutron sources that provide neutron radiation of varying intensities. Neutron diffraction (elastic scattering) techniques are used for analyzing structures; where inelastic neutron scattering is used in studying atomic vibrations and other excitations.
Scattering of fast neutrons
"Fast neutrons" (see neutron temperature) have a kinetic energy above 1 MeV. They can be scattered by condensed matter—nuclei having kinetic energies far below 1 eV—as a valid experimental approximation of an elastic collision with a particle at rest. With each collision, the fast neutron transfers a significant part of its kinetic energy to the scattering nucleus (condensed matter), the more so the lighter the nucleus. And with each collision, the "fast" neutron is slowed until it reaches thermal equilibrium with the material in which it is scattered.
Neutron moderators are used to produce thermal neutrons, which have kinetic energies below 1 eV (T < 500K). Thermal neutrons are used to maintain a nuclear chain reaction in a nuclear reactor, and as a research tool in neutron scattering experiments and other applications of neutron science (see below). The remainder of this article concentrates on the scattering of thermal neutrons.
Neutron-matter interaction
Because neutrons are electrically neutral, they penetrate more deeply into matter than electrically charged particles of comparable kinetic energy, and thus are valuable as probes of bulk properties.
Neutrons interact with atomic nuclei and with magnetic fields from unpaired electrons, causing pronounced interference and energy transfer effects in neutron scattering experiments. Unlike an x-ray photon with a similar wavelength, which interacts with the electron cloud surrounding the nucleus, neutrons interact primarily with the nucleus itself, as described by Fermi's pseudopotential. Neutron scattering and absorption cross sections vary widely from isotope to isotope.
Neutron scattering can be incoherent or coherent, also depending on isotope. Among all isotopes, hydrogen has the highest scattering cross section. Important elements like carbon and oxygen are quite visible in neutron scattering—this is in marked contrast to X-ray scattering where cross sections systematically increase with atomic number. Thus neutrons can be used to analyze materials with low atomic numbers, including proteins and surfactants. This can be done at synchrotron sources but very high intensities are needed, which may cause the structures to change. The nucleus provides a very short range, as isotropic potential varies randomly from isotope to isotope, which makes it possible to tune the (scattering) contrast to suit the experiment.
Scattering almost always presents both elastic and inelastic components. The fraction of elastic scattering is determined by the Debye-Waller factor or the Mössbauer-Lamb factor. Depending on the research question, most measurements concentrate on either elastic or inelastic scattering.
Achieving a precise velocity, i.e. a precise energy and de Broglie wavelength, of a neutron beam is important. Such single-energy beams are termed 'monochromatic', and monochromaticity is achieved either with a crystal monochromator or with a time of flight (TOF) spectrometer. In the time-of-flight technique, neutrons are sent through a sequence of two rotating slits such that only neutrons of a particular velocity are selected. Spallation sources have been developed that can create a rapid pulse of neutrons. The pulse contains neutrons of many different velocities or de Broglie wavelengths, but separate velocities of the scattered neutrons can be determined afterwards by measuring the time of flight of the neutrons between the sample and neutron detector.
Magnetic scattering
The neutron has a net electric charge of zero, but has a significant magnetic moment, although only about 0.1% of that of the electron. Nevertheless, it is large enough to scatter from local magnetic fields inside condensed matter, providing a weakly interacting and hence penetrating probe of ordered magnetic structures and electron spin fluctuations.
Inelastic neutron scattering
Inelastic neutron scattering is an experimental technique commonly used in condensed matter research to study atomic and molecular motion as well as magnetic and crystal field excitations. It distinguishes itself from other neutron scattering techniques by resolving the change in kinetic energy that occurs when the collision between neutrons and the sample is an inelastic one. Results are generally communicated as the dynamic structure factor (also called inelastic scattering law) , sometimes also as the dynamic susceptibility where the scattering vector is the difference between incoming and outgoing wave vector, and is the energy change experienced by the sample (negative that of the scattered neutron). When results are plotted as function of , they can often be interpreted in the same way as spectra obtained by conventional spectroscopic techniques; insofar as inelastic neutron scattering can be seen as a special spectroscopy.
Inelastic scattering experiments normally require a monochromatization of the incident or outgoing beam and an energy analysis of the scattered neutrons. This can be done either through time-of-flight techniques (neutron time-of-flight scattering) or through Bragg reflection from single crystals (neutron triple-axis spectroscopy, neutron backscattering). Monochromatization is not needed in echo techniques (neutron spin echo, neutron resonance spin echo), which use the quantum mechanical phase of the neutrons in addition to their amplitudes.
History
The first neutron diffraction experiments were performed in the 1930s. However it was not until around 1945, with the advent of nuclear reactors, that high neutron fluxes became possible, leading to the possibility of in-depth structure investigations. The first neutron-scattering instruments were installed in beam tubes at multi-purpose research reactors. In the 1960s, high-flux reactors were built that were optimized for beam-tube experiments. The development culminated in the high-flux reactor of the Institut Laue-Langevin (in operation since 1972) that achieved the highest neutron flux to this date. Besides a few high-flux sources, there were some twenty medium-flux reactor sources at universities and other research institutes. Starting in the 1980s, many of these medium-flux sources were shut down, and research concentrated at a few world-leading high-flux sources.
Facilities
Today, most neutron scattering experiments are performed by research scientists who apply for beamtime at neutron sources through a formal proposal procedure. Because of the low count rates involved in neutron scattering experiments, relatively long periods of beam time (on the order of days) are usually required for usable data sets. Proposals are assessed for feasibility and scientific interest.
Techniques
Neutron diffraction
Small angle neutron scattering
Spin Echo Small angle neutron scattering
Neutron reflectometry
Inelastic neutron scattering
Neutron triple-axis spectrometry
Neutron time-of-flight scattering
Neutron backscattering
Neutron spin echo
See also
Neutron transport
LARMOR neutron microscope
Born approximation
References
External links
Free, EU-sponsored e-learning resource for neutron scattering
Neutron scattering - a case study
Neutron Scattering - A primer (LANL-hosted black-and-white version) - An introductory article written by Roger Pynn (Los Alamos National Laboratory)
Podcast Interview with two ILL scientists about neutron science/scattering at the ILL
YouTube video explaining the activities of the Jülich Centre for Neutron Scattering
Neutronsources.org
Science and Innovation with Neutrons in Europe in 2020 (SINE2020)
IAEA neutron beam instrument database
Crystallography
scattering
Neutron
Scattering
de:Neutronenstreuung | Neutron scattering | [
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1,042,263 | https://en.wikipedia.org/wiki/Poynting%27s%20theorem | In electrodynamics, Poynting's theorem is a statement of conservation of energy for electromagnetic fields developed by British physicist John Henry Poynting. It states that in a given volume, the stored energy changes at a rate given by the work done on the charges within the volume, minus the rate at which energy leaves the volume. It is only strictly true in media which is not dispersive, but can be extended for the dispersive case.
The theorem is analogous to the work-energy theorem in classical mechanics, and mathematically similar to the continuity equation.
Definition
Poynting's theorem states that the rate of energy transfer per unit volume from a region of space equals the rate of work done on the charge distribution in the region, plus the energy flux leaving that region.
Mathematically:
where:
is the rate of change of the energy density in the volume.
∇•S is the energy flow out of the volume, given by the divergence of the Poynting vector S.
J•E is the rate at which the fields do work on charges in the volume (J is the current density corresponding to the motion of charge, E is the electric field, and • is the dot product).
Integral form
Using the divergence theorem, Poynting's theorem can also be written in integral form:
where
S is the energy flow, given by the Poynting Vector.
is the energy density in the volume.
is the boundary of the volume. The shape of the volume is arbitrary but fixed for the calculation.
Continuity equation analog
In an electrical engineering context the theorem is sometimes written with the energy density term u expanded as shown. This form resembles the continuity equation:
,
where
ε0 is the vacuum permittivity and μ0 is the vacuum permeability.
is the density of reactive power driving the build-up of electric field,
is the density of reactive power driving the build-up of magnetic field, and
is the density of electric power dissipated by the Lorentz force acting on charge carriers.
Derivation
The rate of work done by the electromagnetic field on the infintesimal charge is given by the Lorentz Force Law as:
(the dot product because from the definition of cross product the cross product of v and B is perpendicular to v.
Where is the volume charge density and is the current density at the point and time where is the velocity of the charge dq.
The rate of work done on the whole charges in the volume V will be the volume integral
By Ampère's circuital law:
(Note that the H and D forms of the magnetic and electric fields are used here. The B and E forms could also be used in an equivalent derivation.)
Substituting this into the expression for rate of work gives:
Using the vector identity :
By Faraday's Law:
giving:
Continuing the derivation requires the following assumptions:
the charges are moving in a medium which is not dispersive.
the total electromagnetic energy density, even for time-varying fields, is given by
It can be shown that:
and
and so:
Returning to the equation for rate of work,
Since the volume is arbitrary, this can be cast in differential form as:
where is the Poynting vector.
Poynting vector in macroscopic media
In a macroscopic medium, electromagnetic effects are described by spatially averaged (macroscopic) fields. The Poynting vector in a macroscopic medium can be defined self-consistently with microscopic theory, in such a way that the spatially averaged microscopic Poynting vector is exactly predicted by a macroscopic formalism. This result is strictly valid in the limit of low-loss and allows for the unambiguous identification of the Poynting vector form in macroscopic electrodynamics.
Alternative forms
It is possible to derive alternative versions of Poynting's theorem. Instead of the flux vector as above, it is possible to follow the same style of derivation, but instead choose , the Minkowski form , or perhaps . Each choice represents the response of the propagation medium in its own way: the form above has the property that the response happens only due to electric currents, while the form uses only (fictitious) magnetic monopole currents. The other two forms (Abraham and Minkowski) use complementary combinations of electric and magnetic currents to represent the polarization and magnetization responses of the medium.
Modification
The derivation of the statement is dependent on the assumption that the materials the equation models can be described by a set of susceptibility properties that are linear, isotropic, homogenous and independent of frequency. The assumption that the materials have no absorption must also be made. A modification to Poynting's theorem to account for variations includes a term for the rate of non-Ohmic absorption in a material, which can be calculated by a simplified approximation based on the Drude model.
Complex Poynting vector theorem
This form of the theorem is useful in Antenna theory, where one has often to consider harmonic fields propagating in the space.
In this case, using phasor notation, and .
Then the following mathematical identity holds:
where is the current density.
Note that in free space, and are real, thus,
taking the real part of the above formula, it expresses the fact that the averaged radiated power flowing through is equal to the work on the charges.
References
External links
Eric W. Weisstein "Poynting Theorem" From ScienceWorld – A Wolfram Web Resource.
Electrodynamics
Eponymous theorems of physics
Circuit theorems | Poynting's theorem | [
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1,042,498 | https://en.wikipedia.org/wiki/Spontaneous%20parametric%20down-conversion | Spontaneous parametric down-conversion (also known as SPDC, parametric fluorescence or parametric scattering) is a nonlinear instant optical process that converts one photon of higher energy (namely, a pump photon) into a pair of photons (namely, a signal photon, and an idler photon) of lower energy, in accordance with the law of conservation of energy and law of conservation of momentum. It is an important process in quantum optics, for the generation of entangled photon pairs, and of single photons.
Basic process
A nonlinear crystal is used to produce pairs of photons from a photon beam. In accordance with the law of conservation of energy and law of conservation of momentum, the pairs have combined energies and momenta equal to the energy and momentum of the original photon. Because the index of refraction changes with frequency (dispersion), only certain triplets of frequencies will be phase-matched so that simultaneous energy and momentum conservation can be achieved. Phase-matching is most commonly achieved using birefringent nonlinear materials, whose index of refraction changes with polarization. As a result of this, different types of SPDC are categorized by the polarizations of the input photon (the pump) and the two output photons (signal and idler). If the signal and idler photons share the same polarization with each other and with the destroyed pump photon it is deemed Type-0 SPDC; if the signal and idler photons share the same polarization to each other, but are orthogonal to the pump polarization, it is Type-I SPDC; and if the signal and idler photons have perpendicular polarizations, it is deemed Type II SPDC.
The conversion efficiency of SPDC is typically very low, with the highest efficiency obtained on the order of 4x10−6 incoming photons for PPLN in waveguides. However, if one half of the pair is detected at any time then its partner is known to be present. The degenerate portion of the output of a Type I down converter is a squeezed vacuum that contains only even photon number terms. The nondegenerate output of the Type II down converter is a two-mode squeezed vacuum.
Example
In a commonly used SPDC apparatus design, a strong laser beam, termed the "pump" beam, is directed at a BBO (beta-barium borate) or lithium niobate crystal. Most of the photons continue straight through the crystal. However, occasionally, some of the photons undergo spontaneous down-conversion with Type II polarization correlation, and the resultant correlated photon pairs have trajectories that are constrained along the sides of two cones whose axes are symmetrically arranged relative to the pump beam. Due to the conservation of momentum, the two photons are always symmetrically located on the sides of the cones, relative to the pump beam. In particular, the trajectories of a small proportion of photon pairs will lie simultaneously on the two lines where the surfaces of the two cones intersect. This results in entanglement of the polarizations of the pairs of photons emerging on those two lines. The photon pairs are in an equal weight quantum superposition of the unentangled states and , corresponding to polarizations of left-hand side photon, right-hand side photon.
Another crystal is KDP (potassium dihydrogen phosphate) which is mostly used in Type I down conversion, where both photons have the same polarization.
Some of the characteristics of effective parametric down-converting nonlinear crystals include:
Nonlinearity: The refractive index of the crystal changes with the intensity of the incident light. This is known as the nonlinear optical response.
Periodicity: The crystal has a regular, repeating structure. This is known as the lattice structure, which is responsible for the regular arrangement of the atoms in the crystal.
Optical anisotropy: The crystal has different refractive indices along different crystallographic axes.
Temperature and pressure sensitivity: The nonlinearity of the crystal can change with temperature and pressure, and thus the crystal should be kept in a stable temperature and pressure environment.
High nonlinear coefficient: Large nonlinear coefficient is desirable, this allow to generate a high number of entangled photons.
High optical damage threshold: Crystal with high optical damage threshold can endure high intensity of the pumping beam.
Transparency in the desired wavelength range: It is important for the crystal to be transparent in the wavelength range of the pump beam for efficient nonlinear interactions
High optical quality and low absorption: The crystal should be high optical quality and low absorption to minimize loss of the pump beam and the generated entangled photons.
History
SPDC was demonstrated as early as 1967 by S. E. Harris, M. K. Oshman, and R. L. Byer, as well as by D. Magde and H. Mahr. It was first applied to experiments related to coherence by two independent pairs of researchers in the late 1980s: Carroll Alley and Yanhua Shih, and Rupamanjari Ghosh and Leonard Mandel. The duality between incoherent (Van Cittert–Zernike theorem) and biphoton emissions was found.
Applications
SPDC allows for the creation of optical fields containing (to a good approximation) a single photon. As of 2005, this is the predominant mechanism for an experimenter to create single photons (also known as Fock states). The single photons as well as the photon pairs are often used in quantum information experiments and applications like quantum cryptography and Bell test experiments.
SPDC is widely used to create pairs of entangled photons with a high degree of spatial correlation. Such pairs are used in ghost imaging, in which information is combined from two light detectors: a conventional, multi-pixel detector that does not view the object, and a single-pixel (bucket) detector that does view the object.
Alternatives
The newly observed effect of two-photon emission from electrically driven semiconductors has been proposed as a basis for more efficient sources of entangled photon pairs. Other than SPDC-generated photon pairs, the photons of a semiconductor-emitted pair usually are not identical but have different energies. Until recently, within the constraints of quantum uncertainty, the pair of emitted photons were assumed to be co-located: they are born from the same location. However, a new nonlocalized mechanism for the production of correlated photon pairs in SPDC has highlighted that occasionally the individual photons that constitute the pair can be emitted from spatially separated points.
See also
Photon upconversion
References
Quantum optics
Articles containing video clips
Light | Spontaneous parametric down-conversion | [
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7,381,179 | https://en.wikipedia.org/wiki/Avrami%20equation | The Avrami equation describes how solids transform from one phase to another at constant temperature. It can specifically describe the kinetics of crystallisation, can be applied generally to other changes of phase in materials, like chemical reaction rates, and can even be meaningful in analyses of ecological systems.
The equation is also known as the Johnson–Mehl–Avrami–Kolmogorov (JMAK) equation. The equation was first derived by Johnson, Mehl, Avrami and Kolmogorov (in Russian) in a series of articles published in the Journal of Chemical Physics between 1939 and 1941. Moreover, Kolmogorov treated statistically the crystallization of a solid in 1937 (in Russian, Kolmogorov, A. N., Izv. Akad. Nauk. SSSR., 1937, 3, 355).
Transformation kinetics
Transformations are often seen to follow a characteristic s-shaped, or sigmoidal, profile where the transformation rates are low at the beginning and the end of the transformation but rapid in between.
The initial slow rate can be attributed to the time required for a significant number of nuclei of the new phase to form and begin growing. During the intermediate period the transformation is rapid as the nuclei grow into particles and consume the old phase while nuclei continue to form in the remaining parent phase.
Once the transformation approaches completion, there remains little untransformed material for further nucleation, and the production of new particles begins to slow. Additionally, the previously formed particles begin to touch one another, forming a boundary where growth stops.
Derivation
The simplest derivation of the Avrami equation makes a number of significant assumptions and simplifications:
Nucleation occurs randomly and homogeneously over the entire untransformed portion of the material.
The growth rate does not depend on the extent of transformation.
Growth occurs at the same rate in all directions.
If these conditions are met, then a transformation of into will proceed by the nucleation of new particles at a rate per unit volume, which grow at a rate into spherical particles and only stop growing when they impinge upon each other. During a time interval , nucleation and growth can only take place in untransformed material. However, the problem is more easily solved by applying the concept of an extended volume – the volume of the new phase that would form if the entire sample was still untransformed. During the time interval to the number of nuclei N that appear in a sample of volume V will be given by
where is one of two parameters in this simple model: the nucleation rate per unit volume, which is assumed to be constant. Since growth is isotropic, constant and unhindered by previously transformed material, each nucleus will grow into a sphere of radius , and so the extended volume of due to nuclei appearing in the time interval will be
where is the second of the two parameters in this simple model: the growth velocity of a crystal, which is also assumed constant. The integration of this equation between and will yield the total extended volume that appears in the time interval:
Only a fraction of this extended volume is real; some portion of it lies on previously transformed material and is virtual. Since nucleation occurs randomly, the fraction of the extended volume that forms during each time increment that is real will be proportional to the volume fraction of untransformed . Thus
rearranged
and upon integration:
where Y is the volume fraction of ().
Given the previous equations, this can be reduced to the more familiar form of the Avrami (JMAK) equation, which gives the fraction of transformed material after a hold time at a given temperature:
where , and .
This can be rewritten as
which allows the determination of the constants n and from a plot of vs . If the transformation follows the Avrami equation, this yields a straight line with slope n and intercept .
Final crystallite (domain) size
Crystallization is largely over when reaches values close to 1, which will be at a crystallization time defined by , as then the exponential term in the above expression for will be small. Thus crystallization takes a time of order
i.e., crystallization takes a time that decreases as one over the one-quarter power of the nucleation rate per unit volume, , and one over the three-quarters power of the growth velocity . Typical crystallites grow for some fraction of the crystallization time and so have a linear dimension , or
i.e., the one quarter power of the ratio of the growth velocity to the nucleation rate per unit volume. Thus the size of the final crystals only depends on this ratio, within this model, and as we should expect, fast growth rates and slow nucleation rates result in large crystals. The average volume of the crystallites is of order this typical linear size cubed.
This all assumes an exponent of , which is appropriate for the uniform (homogeneous) nucleation in three dimensions. Thin films, for example, may be effectively two-dimensional, in which case if nucleation is again uniform the exponent . In general, for uniform nucleation and growth, , where is the dimensionality of space in which crystallization occurs.
Interpretation of Avrami constants
Originally, n was held to have an integer value between 1 and 4, which reflected the nature of the transformation in question. In the derivation above, for example, the value of 4 can be said to have contributions from three dimensions of growth and one representing a constant nucleation rate. Alternative derivations exist, where n has a different value.
If the nuclei are preformed, and so all present from the beginning, the transformation is only due to the 3-dimensional growth of the nuclei, and n has a value of 3.
An interesting condition occurs when nucleation occurs on specific sites (such as grain boundaries or impurities) that rapidly saturate soon after the transformation begins. Initially, nucleation may be random, and growth unhindered, leading to high values for n (3 or 4). Once the nucleation sites are consumed, the formation of new particles will cease.
Furthermore, if the distribution of nucleation sites is non-random, then the growth may be restricted to 1 or 2 dimensions. Site saturation may lead to n values of 1, 2 or 3 for surface, edge and point sites respectively.
Applications in biophysics
The Avrami equation was applied in cancer biophysics in two aspects. First aspect is connected with tumor growth and cancer cells kinetics, which can be described by the sigmoidal curve. In this context the Avrami function was discussed as an alternative to the widely used Gompertz curve. In the second aspect the Avrami nucleation and growth theory was used together with multi-hit theory of carcinogenesis to show how the cancer cell is created. The number of oncogenic mutations in cellular DNA can be treated as nucleation particles which can transform whole DNA molecule into cancerous one (neoplastic transformation). This model was applied to clinical data of gastric cancer, and shows that Avrami's constant n is between 4 and 5 which suggest the fractal geometry of carcinogenic dynamics. Similar findings were published for breast and ovarian cancers, where n=5.3.
Multiple Fitting of a Single Dataset (MFSDS)
The Avrami equation was used by Ivanov et al. to fit multiple times a dataset generated by another model, the so called αDg to а sequence of the upper values of α, always starting from α=0, in order to generate a sequence of values of the Avrami parameter n. This approach was shown effective for a given experimental dataset, see the plot, and the n values obtained follow the general direction predicted by fitting multiple times the α21 model.
References
External links
IUPAC Compendium of Chemical Terminology 2nd ed. (the "Gold Book"), Oxford (1997)
Crystallography
Equations | Avrami equation | [
"Physics",
"Chemistry",
"Materials_science",
"Mathematics",
"Engineering"
] | 1,656 | [
"Mathematical objects",
"Materials science",
"Equations",
"Crystallography",
"Condensed matter physics"
] |
7,381,751 | https://en.wikipedia.org/wiki/Permissiveness%20%28endocrinology%29 | In endocrinology, permissiveness is a biochemical phenomenon in which the presence of one hormone is required in order for another hormone to exert its full effects on a target cell. Hormones can interact in permissive, synergistic, or antagonistic ways. The chemical classes of hormones include amines, polypeptides, glycoproteins and steroids. Permissive hormones act as precursors to active hormones and may be classified as either prohormones or prehormones. It stimulate the formation of receptors of that hormone.
Examples
Thyroid hormone increases the number of beta-adrenergic receptors available for epinephrine at the latter's target cell, thereby increasing epinephrine's effect on that cell. Specially in cardiac cell. Without the thyroid hormone, epinephrine would have only a weak effect.
Cortisol is required for the response of vascular and bronchial smooth muscle to catecholamines. Cortisol is also required for the lipolytic effect of catecholamines, ACTH, and growth hormone on fat cells. Cortisol is also required for the calorigenic effects of glucagon and catecholamines.
The effects of a hormone in the body depend on its concentration. Permissive actions of glucocorticoids like cortisol generally occur at low concentrations. Abnormally high amounts of a hormone can result in atypical effects. Glucocorticoids function by attaching to cytoplasmic receptors to either enhance or suppress changes in the transcription of DNA and thus the synthesis of proteins. Glucocorticoids also inhibit the secretion of cytokines via post-translational modification effects.
References
Biology terminology | Permissiveness (endocrinology) | [
"Chemistry",
"Biology"
] | 364 | [
"Biochemistry stubs",
"Biotechnology stubs",
"Biochemistry",
"nan"
] |
7,382,202 | https://en.wikipedia.org/wiki/IRE%20%28unit%29 | The IRE unit is used in the measurement of composite video signals. Its name is derived from the initials of the Institute of Radio Engineers.
A value of 100 IRE is defined to be +714 mV in an analog NTSC video signal. A value of 0 IRE corresponds to the voltage value of 0 mV, the signal value during the blanking period. The sync pulse is normally 40 IRE below this 0 IRE value, so the total range covered from peak to trough of an all white signal would be 140 IRE.
Video signals use the "IRE" unit instead of DC voltages to describe levels and amplitudes. Based on a standard 1 Vpp NTSC composite-video signal that swings from -286 mV (sync tip) to +714 mV (peak video), a 140 IRE peak-to-peak convention is established. Thus, one NTSC IRE unit is 7.143 mV ( V or mV), where -40 IRE is equivalent to -285.7 mV, and +100 IRE is equivalent to +714.3 mV. 0 IRE is equivalent to 0 V. The black level is equivalent to 53.57 mV (7.5 IRE).
The PAL video signal is slightly different in that it swings from -300 mV to +700 mV, instead. Thus, one PAL IRE unit is 7 mV, where -43 IRE is equivalent to -300 mV at the sync tip, and +100 IRE is equivalent to +700 mV at the peak video level. Black level is the same as the blanking level 0 mV (0 IRE).
The reason IRE is a relative measurement (percent) is because a video signal may be any amplitude. This unit is used in the ITU recommendations BT.470 and BT.1700 which define PAL, NTSC and SECAM:
References
Units of measurement
Broadcast engineering
Video formats
Television technology
ITU-R recommendations | IRE (unit) | [
"Mathematics",
"Technology",
"Engineering"
] | 399 | [
"Information and communications technology",
"Broadcast engineering",
"Television technology",
"Quantity",
"Electronic engineering",
"Units of measurement"
] |
7,383,487 | https://en.wikipedia.org/wiki/Oospore | An oospore is a thick-walled sexual spore that develops from a fertilized oosphere in some algae, fungi, and oomycetes. They are believed to have evolved either through the fusion of two species or the chemically induced stimulation of mycelia, leading to oospore formation.
In Oomycetes, oospores can also result from asexual reproduction, by apomixis. These haploid, non-motile spores are the site of meiosis and karyogamy in oomycetes.
A dormant oospore, when observed under an electron microscope, has led researchers to draw conclusion that there is only a single central globule with other storage bodies surrounding it.
References
Reproduction | Oospore | [
"Biology"
] | 157 | [
"Biological interactions",
"Behavior",
"Reproduction"
] |
7,383,878 | https://en.wikipedia.org/wiki/Overdrafting | Overdrafting is the process of extracting groundwater beyond the equilibrium yield of an aquifer. Groundwater is one of the largest sources of fresh water and is found underground. The primary cause of groundwater depletion is the excessive pumping of groundwater up from underground aquifers. Insufficient recharge can lead to depletion, reducing the usefulness of the aquifer for humans. Depletion can also have impacts on the environment around the aquifer, such as soil compression and land subsidence, local climatic change, soil chemistry changes, and other deterioration of the local environment.
There are two sets of yields: safe yield and sustainable yield. Safe yield is the amount of groundwater that can be withdrawn over a period of time without exceeding the long-term recharge rate or affecting the aquifer integrity. Sustainable yield is the amount of water extraction that can be sustained indefinitely without negative hydrological impacts, taking into account both recharge rate and surface water impacts.
There are two types of aquifers: confined and unconfined. In confined aquifers, there is an overbearing layer called an aquitard, which contains impermeable materials through which groundwater cannot be extracted. In unconfined aquifers, there is no aquitard, and groundwater can be freely extracted from the surface. Extracting groundwater from unconfined aquifers is like borrowing the water: it has to be recharged at a proper rate. Recharge can happen through artificial recharge and natural recharge.
Mechanism
When groundwater is extracted from an aquifer, a cone of depression is created around the well. As the drafting of water continues, the cone increases in radius. Extracting too much water (overdrafting) can lead to negative impacts such as a drop of the water table, land subsidence, and loss of surface water reaching the streams. In extreme cases, the supply of water that naturally recharges the aquifer is pulled directly from streams and rivers, lowering their water levels. This affects wildlife, as well as humans who might be using the water for other purposes.
The natural process of aquifer recharge takes place through the percolation of surface water. An aquifer may be artificially recharged, such as by pumping reclaimed water from wastewater management projects directly into the aquifer. An example of is the Orange County Water District in California. This organization takes wastewater, treats it to a proper level, and then systematically pumps it back into the aquifers for artificial recharge.
Since every groundwater basin recharges at a different rate depending on precipitation, vegetative cover, and soil conservation practices, the quantity of groundwater that can be safely pumped varies greatly among regions of the world and even within provinces. Some aquifers require a very long time to recharge, and thus overdrafting can effectively dry up certain sub-surface water supplies. Subsidence occurs when excessive groundwater is extracted from rocks that support more weight when saturated. This can lead to a capacity reduction in the aquifer.
Changes in freshwater availability stem from natural and human activities (in conjunction with climate change) that interfere with groundwater recharge patterns. One of the leading anthropogenic activities causing groundwater depletion is irrigation. Roughly 40% of global irrigation is supported by groundwater, and irrigation is the primary activity causing groundwater storage loss across the U.S.
Around the world
This ranking is based on the amount of groundwater each country uses for agriculture. This issue is becoming significant in the United States (most notably in California), but it has been an ongoing problem in other parts of the world, such as was documented in Punjab, India, in 1987.
United States
In the U.S., an estimated 800 km3 of groundwater was depleted during the 20th century. The development of cities and other areas of highly concentrated water usage has created a strain on groundwater resources. In post-development scenarios, interactions between surface water and groundwater are reduced; there is less intermixing between the surface and subsurface (interflow), leading to depleted water tables.
Groundwater recharge rates are also affected by rising temperatures which increase surface evaporation and transpiration, resulting in decreased water content of the soil. Anthropogenic changes to groundwater storage, such as over-pumping and the depletion of water tables combined with climate change, effectively reshape the hydrosphere and impact the ecosystems that depend on the groundwater.
Accelerated decline in subterranean reservoirs
According to a 2013 report by research hydrologist Leonard F. Konikow at the United States Geological Survey (USGS), the depletion of the Ogallala Aquifer between 20012008 is about 32% of the cumulative depletion during the entire 20th century. In the United States, the biggest users of water from aquifers include agricultural irrigation, and oil and coal extraction. According to Konikow, "Cumulative total groundwater depletion in the United States accelerated in the late 1940s and continued at an almost steady linear rate through the end of the century. In addition to widely recognized environmental consequences, groundwater depletion also adversely impacts the long-term sustainability of groundwater supplies to help meet the Nation's water needs."
As reported by another USGS study of withdrawals from 66 major US aquifers, the three greatest uses of water extracted from aquifers were irrigation (68%), public water supply (19%), and "self-supplied industrial" (4%). The remaining 8% of groundwater withdrawals were for "self-supplied domestic, aquaculture, livestock, mining, and thermoelectric power uses."
Environmental impacts
Groundwater extraction for use in water supplies lowers the overall water table, the level that groundwater sits at in an area. The lowering water table can diminish streamflow and reduce water level in other water bodies such as wetlands and lakes. In Karst systems, large-scale groundwater withdrawal can lead to sinkholes or groundwater-related subsidence. The overdrafting leads to the pressure in limestone containments to become unstable and sediments to collapse, creating a sinkhole. Overdrafting in coastal regions can lead to the reduction of water pressure in an aquifer, allowing saltwater intrusion. If saltwater contaminates a freshwater aquifer, that aquifer can no longer be used as a reliable source of freshwater for settlements and cities. Artificial recharge may return fresh water pressure to halt saltwater intrusion. However, this method can be economically inefficient and unavailable due to the high cost of the process.
When aquifers or groundwater wells experience overdraft, chemical concentrations in the water may change. Chemicals such as calcium, magnesium, sodium, carbonate, bicarbonate, chloride, and sulfate can be found in groundwater sources. Changes to water quality as a result of overdrafting may make it unsafe for human consumption; rendering the groundwater sources unusable as a source of drinking water.
Overdrafting can also affect organisms living within groundwater aquifers known as stygobionts Loss of habitat for these creatures through overdrafting has reduced biodiversity in certain areas.
Environmental impacts of overdrafting include:
Groundwater-related subsidence: the collapse of land due to lack of support (from the water that is being depleted). The first recorded case of land subsidence was in the 1940s. Land subsidence can be as little as local land collapsing or as large as an entire region's land being lowered. The subsidence can lead to infrastructural and ecosystem damage.
Lowering of the water table, which makes water harder to reach streams and rivers
Reduction of water volume in streams and lakes because their supply of water is being diminished by surface water recharging the aquifers
Impacts on animals that depend on streams and lakes for food, water, and habitat
Deterioration to water quality
Increase in the cost of water to the consumer due to a lower water table—more energy is needed to pump from a greater depth, so operating costs increase for companies, who pass on the expense to the consumer
Decrease in crop production from lack of water
Disturbances to the water cycle
Groundwater related subsidence
Socio-economic effects
Overdrafting has socio-economic impacts due to cost inequities that increase as the water table drops. As the water table drops, deeper wells are required to reach water in the aquifer. This not only requires deepening of already existing wells, but also digging new wells. Both processes are expensive. Research from Punjab found that the high cost of technology to continue water access hurts small landholders more than it does large landholders because large landholders have more resources "to invest in technology." Therefore, small landholders, who traditionally have a lower income than large landholders, are unable to benefit from the technology that allows greater water access. This creates a cycle of inequity as small landholders that are dependent on agriculture have less water to irrigate their land, producing a lower output of crops.
Additionally, overdrafting has socio-economic impacts due to prior appropriation laws. Prior appropriation rights declare that the first person to use water from a water source will maintain the right to water. These rights result in socio-economic inequities as businesses and/or larger landholders who have a higher income can maintain their water rights. Meanwhile, new businesses or smaller landholders have less access to water, resulting in less ability to profit. Due to this inequity, small farmers in Punjab with less water rights tend to grow maize or less productive rice; meanwhile, larger landholders in Punjab can use more land for rice because they have access to water.
Possible solutions
Artificial Recharge:
Since recharge is the natural replenishment of water, artificial recharge is the man-made replenishment of groundwater, though there is only a limited amount of suitable water available for replenishing.
Water Conservation Techniques:
Other solutions include implementing water conservation techniques to decrease overdrafting. These include improving governance to ensure proper water management, incentivizing water conservation, improving agriculture techniques to ensure water use is efficient, changing diets to crops that require less water, and investing in infrastructure that uses water sustainably. The state of California has implemented some water conservation techniques due to droughts in the state. Some of their techniques include prohibitions on: 1) outdoor watering that runs onto sidewalks or other on hard surfaces that don't absorb water, 2) washing vehicles with a hose that does not have a shutoff handle, 3) watering within 48 hours after a quarter inch of rain, and 4) watering commercial/industrial decorative grass.
Water Conservation Incentivization:
Techniques used by California in emergency situations are useful; however, incentive to follow through on these is important. The city of Spokane has a program to incentivize sustainable landscapes called SpokaneScape. This program incentivizes water efficient landscapes by offering homeowners up to $500 in credit on their utility bill if they adapt their yards to water efficient plants.
See also
Cone of depression
Groundwater recharge
Groundwater-related subsidence
Drinking water
Overexploitation
Water crisis
Human overpopulation
References
External links
The Perils of Groundwater Pumping, Issues in Science and Technology
Aquifers
Environmental impact of agriculture
Environmental issues with water
Water supply
Water and the environment | Overdrafting | [
"Chemistry",
"Engineering",
"Environmental_science"
] | 2,345 | [
"Hydrology",
"Aquifers",
"Water supply",
"Environmental engineering"
] |
7,384,821 | https://en.wikipedia.org/wiki/Guerbet%20reaction | The Guerbet reaction, named after Marcel Guerbet (1861–1938), is an organic reaction that converts a primary alcohol into its β-alkylated dimer alcohol with loss of one equivalent of water. The process is of interest because it converts simple inexpensive feedstocks into more valuable products. Its main disadvantage is that the reaction produces mixtures.
Scope and applications
The original 1899 publication concerned the conversion of n-butanol to 2-ethylhexanol. 2-ethylhexanol is however more easily prepared by alternative methods (from butyraldehyde by aldol condensation).
Instead, the Guerbet reaction is mainly applied to fatty alcohols to afford oily products, which are called Guerbet alcohols. They are of commercial interest to as components of cosmetics, plasticizers, and related applications. The reaction is conducted in the temperature range 180-360 °C, often in a sealed reactor. The reaction requires alkali metal hydroxides or alkoxides. Catalysts such as Raney Nickel are required to facilitate the hydrogen transfer steps.
While the Guerbet reaction is traditionally (and commercially) focused on fatty alcohols, it has been investigated for the dimerization of ethanol to butanol.
Organometallic catalysts have been investigated. A small amount of the diene 1,7-octadiene is required as a hydrogen acceptor.
Mechanism
The reaction mechanism for this reaction is a four-step sequence. In the first step the alcohol is oxidized to the aldehyde. These intermediates then react in an aldol condensation to the allyl aldehyde which the hydrogenation catalyst then reduces to the alcohol.
The Cannizzaro reaction is a competing reaction when two aldehyde molecules react by disproportionation to form the corresponding alcohol and carboxylic acid. Another side reaction is the Tishchenko reaction.
See also
Oxo alcohols - a different reaction which gives similar products
Guerbet alcohols
2-Ethyl-1-butanol
2-Ethylhexanol
2-Propylheptan-1-ol
2-Butyl-1-octanol
2-Butyl-1-octanol
References
External links
A Review of Guerbet Chemistry Anthony J. O’Lenick, Jr. https://web.archive.org/web/20110209074739/http://www.zenitech.com/ Link
Condensation reactions
Name reactions
Fatty alcohols | Guerbet reaction | [
"Chemistry"
] | 532 | [
"Name reactions",
"Condensation reactions",
"Organic reactions"
] |
7,385,054 | https://en.wikipedia.org/wiki/List%20of%20gemstones%20by%20species | This is a list of gemstones, organized by species and types.
Minerals
There are over 300 types of minerals that have been used as gemstones. Such as:
A–B
Actinolite
Nephrite ()
Adamite
Aegirine
Afghanite
Agrellite
Algodonite
Alunite
Amblygonite
Analcime
Anatase
Andalusite
Chiastolite
Andesine
Anglesite
Anhydrite
Annabergite
Anorthite
Antigorite
Bowenite
Apatite
Apophyllite
Aragonite
Arfvedsonite
Asbestos
Astrophyllite
Atacamite
Augelite
Austinite
Axinite group:
Ferroaxinite
Magnesioaxinite
Manganaxinite
Tinzenite
Azurmalachite
Azurite
Baryte
Bastnaesite
Bayldonite
Benitoite
Beryl subgroup:
Aquamarine
Maxixe
Emerald
Trapiche emerald ()
Goshenite
Golden beryl
Heliodor
Morganite
Red beryl (Bixbite)
Beryllonite
Beudantite
Bismutotantalite
Biotite
Boleite
Boracite
Bornite
Brazilianite
Breithauptite
Brookite
Brucite
Bustamite
Bytownite
C–F
Calcite
Manganoan calcite ()
Caledonite
Canasite
Cancrinite
Vishnevite
Carletonite
Carnallite
Cassiterite
Catapleiite
Cavansite
Celestite
Ceruleite
Cerussite
Chabazite
Chalcopyrite
Chambersite
Charlesite
Charoite
Childrenite
Chiolite
Chondrodite
Chrysoberyl
Alexandrite ()
Cymophane
Chromite
Chrysocolla
Chrysotile
Cinnabar
Clinochlore
Clinohumite
Clinozoisite
Clintonite
Cobaltite
Colemanite
Cordierite
Iolite ()
Cornwallite
Corundum
Ruby ()
Sapphire ()
Padparadscha
Golden sheen sapphire
Covellite
Creedite
Crocoite
Cryolite
Cumberlandite
Cuprite
Danburite
Datolite
Descloizite
Diamond
Bort
Ballas
Diaspore
Dickinsonite
Diopside
Dioptase
Dolomite
Dumortierite
Ekanite
Elbaite
Enstatite
Bronzite
Hypersthene
Eosphorite
Epidote
Piemontite
Erythrite
Esperite
Ettringite
Euclase
Eudialyte
Euxenite
Fayalite
Feldspar subgroup:
Andesine
Albite
Anorthite
Anorthoclase
Amazonite
Bytownite
Celsian
Microcline
Moonstone
Adularia ()
Rainbow ()
Orthoclase
Unakite
Plagioclase
Albite
Labradorite
Oligoclase
Sanidine
Sunstone
Oregon sunstone
Rainbow lattice sunstone
Fergusonite
Ferroaxinite
Fluorapatite
Fluorapophyllite
Fluorite
Forsterite
Friedelite
G–L
Gadolinite
Gahnite
Gahnospinel
Garnet group:
Pyralspite
Almandine
Pyrope
Spessartine
Ugrandite
Andradite
Demantoid
Melanite
Topazolite
Grossular
Hessonite
Hydrogrossular
Tsavorite
Uvarovite
Almandine-pyrope
Rhodolite
Andradite-grossular
Grandite (Mali garnet)
Pyrope-almandine-spessartine
Malaia garnet
Pyrope-spessartine
Umbalite
Gaspeite
Gaylussite
Gibbsite
Glaucophane
Goethite
Goosecreekite
Grandidierite
Gypsum
Gyrolite
Halite
Hambergite
Hanksite
Hardystonite
Hauyne
Helenite
Hematite
Hemimorphite
Herderite
Hexagonite
Hibonite
Hiddenite
Hodgkinsonite
Holtite
Howlite
Huebnerite
Humite
Hureaulite
Hurlbutite
Hyperitdiabas
Ilmenite
Inderite
Jade
Jadeite
Chloromelanite
Nephrite
Jasper
Jeremejevite
Kainite
Kämmererite
Kaolinite
Kornerupine
Kutnohorite
Kurnakovite
Kyanite
Langbeinite
Lawsonite
Lazulite
Lazurite
Legrandite
Lepidolite
Leucite
Leucophanite
Linarite
Lizardite
Londonite
Ludlamite
Ludwigite
M–Q
Magnesite
Malachite
Marialite-meionite
Wernerite ()
Marcasite
Meliphanite
Mellite
Mesolite
Microcline
Microlite
Milarite
Millerite
Mimetite
Monazite
Mordenite
Mottramite
Muscovite
Fuchsite ()
Musgravite
Nambulite
Narsarsukite
Natrolite
Nepheline
Neptunite
Nickeline (Niccolite)
Nosean
Nuummite
Olivine
Opal
Fire opal
Moss opal
Painite
Palygorskite
Papagoite
Pargasite
Parisite
Pectolite
Larimar
Pentlandite
Peridot
Periclase
Perthite
Petalite (castorite)
Pezzottaite
Phenakite
Phlogopite
Phosgenite
Phosphophyllite
Phosphosiderite
Piemontite
Pietersite
Plumbogummite
Pollucite
Polyhalite
Poudretteite
Powellite
Prehnite
Prismatine
Prosopite
Proustite
Psilomelane
Pumpellyite
Chlorastrolite ()
Purpurite
Pyrite
Pyrargyrite
Pyromorphite
Pyrophyllite
Pyroxmangite
Pyrrhotite
Quartz
Amethyst ()
Ametrine ()
Aventurine ()
Chalcedony ()
Agate
Iris agate
Onyx
Sardonyx
Bloodstone (Heliotrope)
Carnelian
Chrome chalcedony
Chrysoprase
Dendritic agate
Moss agate
Fire agate (iridescent )
Jasper
Petrified wood
Sard
Citrine ()
Druzy ()
Flint ()
Herkimer diamond ()
Milky quartz ()
Prasiolite ()
Radiolarite ()
Rose quartz ()
Rock crystal ()
Shocked quartz ()
Smoky quartz ()
Quartzite
R–Z
Realgar
Rhodizite
Rhodochrosite
Rhodonite
Richterite
Riebeckite
Crocidolite ()
Rosasite
Rutile
Samarskite
Sanidine
Sapphirine
Sarcolite
Scapolite
Marialite
Meionite
Scheelite
Schizolite
Scorodite
Selenite
Sellaite
Senarmontite
Sepiolite (Meerschaum)
Sérandite
Seraphinite
Serendibite
Serpentine subgroup
Antigorite
Bowenite
Chrysotile
Lizardite
Stichtite
Shattuckite
Shigaite
Shortite
Shungite
Siderite
Sillimanite
Simpsonite
Sinhalite
Smaltite
Smithsonite
Sodalite
Hackmanite ()
Sogdianite
Sperrylite
Spessartite
Sphalerite
Spinel
Ceylonite ()
Spodumene
Hiddenite ()
Kunzite ()
Triphane ()
Spurrite
Staurolite
Stibiotantalite
Stichtite
Stolzite
Strontianite
Strontium titanate
Sulfur
Sugilite
Bustamite ()
Richterite ()
Sylvite
Taaffeite
Talc
Tantalite
Tektites
Moldavite
Tephroite
Thomsonite
Thaumasite
Tinaksite
Titanite (sphene)
Topaz
Tourmaline subgroup:
Achroite ()
Chrome ()
Dravite
Elbaite
Fluor-liddicoatite
Indicolite
Olenite
Paraiba ()
Rossmanite
Rubellite ()
Tremolite
Hexagonite ()
Triphylite
Triplite
Tugtupite
Turquoise
Ulexite
Ussingite
Vanadinite
Variscite
Väyrynenite
Vesuvianite (idocrase)
Californite ()
Villiaumite
Vivianite
Vlasovite
Wardite
Wavellite
Weloganite
Whewellite
Wilkeite
Willemite
Witherite
Wollastonite
Wulfenite
Wurtzite
Xonotlite
Yugawaralite
Zektzerite
Zeolites
Analcime
Apophyllite
Chabazite
Goosecreekite
Natrolite
Scolecite
Stellerite
Stilbite
Thomsonite
Zincite
Zinnwaldite
Zircon
Jacinth ()
Zoisite
Tanzanite ()
Thulite ()
Zultanite
Zunyite
Artificial and lab created
There are a number of artificial and lab grown minerals used to produce gemstones. These include:
Lab alexandrite
Lab corundum
Cubic zirconia
Lab diamond
Lab emerald
Fordite
Gadolinium gallium garnet
Lab moissanite
Synthetic opal
Metal-coated crystals hyped as rainbow quartz
Lab spinel
Synthetic turquoise
Terbium gallium garnet
Trinitite
Yttrium aluminium garnet
Yttrium iron garnet
Organic
There are a number of organic materials used as gems, including:
Amber
Ammolite
Ammonoidea
Bone
Copal
Coral
Ivory
Jet
Nacre (Mother of pearl)
Operculum
Pearl
Seashell
Rocks
Some rocks are used as gems, including:
Anthracite
Anyolite
Bauxite
Concretions
Bloodstone (Heliotrope)
Eilat stone
Epidosite
Glimmerite
Goldstone (glittering glass)
Tiger's eye
Helenite (artificial glass made from volcanic ash)
Iddingsite
Kimberlite
Lamproite
Lapis lazuli
Libyan desert glass
Llanite
Maw sit sit
Moldavite
Obsidian
Apache tears
Pallasite
Peridotite (also known as olivinite)
Siilinjärvi carbonatite
Soapstone (also known as steatite)
Tactite
Tiger's eye
Unakite
Chatoyant gems
Some minerals made into gemstones may display a chatoyancy or cat's eye effect, these include:
Actinolite
Andalusite
Apatite
Beryl
Aquamarine
Emerald
Heliodor
Morganite
Beryllium
Beryllonite
Calcite
Cerussite
Chrysoberyl
Danburite
Diaspore
Diopside
Enstatite
Garnet
Grandidierite
Hawk's eye
Hypersthene
Iolite
Kornerupine
Kunzite
Kyanite
Moonstone
Opal
Peridot
Peristerite (Albite variety)
Pezzottaite
Phenakite
Prasiolite
Prehnite
Quartz
Rhodonite
Rutile
Scapolite
Selenite
Serpentine
Antigorite
Bowenite
Sillimanite
Smoky Quartz
Spinel
Sunstone
Tanzanite
Tiger's Eye
Topaz
Tourmaline
Ulexite
Zircon
Asterism
Corundum
Ruby
Sapphire
Diopside
See also
List of individual gemstones
References
Further reading
Gemstones of the World revised 5th edition, 2013 by Walter Schumann
Smithsonian Handbook: Gemstones by Cally Hall, 2nd ed. 2002
Species | List of gemstones by species | [
"Physics"
] | 2,270 | [
"Materials",
"Gemstones",
"Matter"
] |
7,385,565 | https://en.wikipedia.org/wiki/Thue%20number | In the mathematical area of graph theory, the Thue number of a graph is a variation of the chromatic index, defined by and named after mathematician Axel Thue, who studied the squarefree words used to define this number.
Alon et al. define a nonrepetitive coloring of a graph to be an assignment of colors to the edges of the graph, such that there does not exist any even-length simple path in the graph in which the colors of the edges in the first half of the path form the same sequence as the colors of the edges in the second half of the path. The Thue number of a graph is the minimum number of colors needed in any nonrepetitive coloring.
Variations on this concept involving vertex colorings or more general walks on a graph have been studied by several authors including Barát and Varjú, Barát and Wood (2005), Brešar and Klavžar (2004), and Kündgen and Pelsmajer.
Example
Consider a pentagon, that is, a cycle of five vertices. If we color the edges with two colors, some two adjacent edges will have the same color x; the path formed by those two edges will have the repetitive color sequence xx. If we color the edges with three colors, one of the three colors will be used only once; the path of four edges formed by the other two colors will either have two consecutive edges or will form the repetitive color sequence xyxy. However, with four colors it is not difficult to avoid all repetitions. Therefore, the Thue number of C5 is four.
Results
Alon et al. use the Lovász local lemma to prove that the Thue number of any graph is at most quadratic in its maximum degree; they provide an example showing that for some graphs this quadratic dependence is necessary. In addition they show that the Thue number of a path of four or more vertices is exactly three, that the Thue number of any cycle is at most four, and that the Thue number of the Petersen graph is exactly five.
The known cycles with Thue number four are C5, C7, C9, C10, C14, and C17. Alon et al. conjecture that the Thue number of any larger cycle is three; they verified computationally that the cycles listed above are the only ones of length ≤ 2001 with Thue number four. Currie resolved this in a 2002 paper, showing that all cycles with 18 or more vertices have Thue number 3.
Computational complexity
Testing whether a coloring has a repetitive path is in NP, so testing whether a coloring is nonrepetitive is in co-NP, and Manin showed that it is co-NP-complete. The problem of finding such a coloring belongs to in the polynomial hierarchy, and again Manin showed that it is complete for this level.
References
External links
Graph invariants
Graph coloring
Combinatorics on words | Thue number | [
"Mathematics"
] | 599 | [
"Graph coloring",
"Graph theory",
"Combinatorics",
"Graph invariants",
"Mathematical relations",
"Combinatorics on words"
] |
4,237,693 | https://en.wikipedia.org/wiki/Spurline | The spurline is a type of radio-frequency and microwave distributed element filter with band-stop (notch) characteristics, most commonly used with microstrip transmission lines. Spurlines usually exhibit moderate to narrow-band rejection, at about 10% around the central frequency.
Spurline filters are very convenient for dense integrated circuits because of their inherently compact design and ease of integration: they occupy surface that corresponds only to a quarter-wavelength transmission line.
Structure description
It consists of a normal microstrip line breaking into a pair of smaller coupled lines that rejoin after a quarter-wavelength distance. Only one of the input ports of the coupled lines is connected to the feed microstrip, as shown in the figure below. The orange area of the illustration is the microstrip transmission line conductor and the gray color the exposed dielectric.
Where is the wavelength corresponding to the central rejection frequency of the bandstop filter, measured - of course - in the microstrip line material. This is the most important parameter of the filter that sets the rejection band.
The distance between the two coupled lines can be selected appropriately to fine-tune the filter. The smaller the distance, the narrower the stop-band in terms of rejection. Of course that is limited by the circuit-board printing resolution, and it is usually considered at about 10% of the input microstrip width.
The gap between the input microstrip line and the one open-circuited line of the coupler has a negligible effect on the frequency response of the filter. Therefore, it is considered approximately equal to the distance of the two coupled lines.
Printed antennae
Spurlines can also be used in printed antennae such as the planar inverted-F antenna. The additional resonances can be designed to widen the antenna bandwidth or to create multiple bands, for instance, for a tri-band mobile phone.
History
A spurline filter was first proposed by Schiffman and Matthaei in stripline form in 1964. Bates adapted the design for microstrip in 1977. Nguyen and Hsieh improved the analysis for microstrip implementations in 1983.
References
C. Nguyen and K. Chang, “On the analysis and design of spurline bandstop filters,” IEEE Trans. Microw. Theory Tech., vol. 33, no. 12, pp. 1416–1421, Dec. 1985.
Primary sources
B. M. Schiffman; G. L. Matthaei, "Exact design of band-stop microwave filters", IEEE Transactions on Microwave Theory and Techniques, vol. 12, iss. 1, pp. 6-15, 1964.
R. N. Bates, "Design of microstrip spur-line band-stop filters", IEEE Journal on Microwave Optics and Acoustics, vol. 1, iss. 6, pp. 209-204, November 1977.
C. Nguyen; C. Hsieh, Millimeter wave printed circuit spurline filters", IEEE MTT-S International Microwave Symposium Digest, pp. 98-100, 1983.
Microwave technology
Distributed element circuits | Spurline | [
"Engineering"
] | 624 | [
"Electronic engineering",
"Distributed element circuits"
] |
4,237,747 | https://en.wikipedia.org/wiki/Simulated%20fluorescence%20process%20algorithm | The Simulated Fluorescence Process (SFP) is a computing algorithm used for scientific visualization of 3D data from, for example, fluorescence microscopes. By modeling a physical light/matter interaction process, an image can be computed which shows the data as it would have appeared in reality when viewed under these conditions.
Principle
The algorithm considers a virtual light source producing excitation light that illuminates the object. This casts shadows either on parts of the object itself or on other objects below it. The interaction between the excitation light and the object provokes the emission light, which also interacts with the object before it finally reaches the eye of the viewer.
See also
Computer graphics lighting
Rendering (computer graphics)
References
External links
Freeware SFP renderer
Computational science
Computer graphics algorithms
Visualization (graphics)
Microscopes
Microscopy
Fluorescence | Simulated fluorescence process algorithm | [
"Physics",
"Chemistry",
"Astronomy",
"Mathematics",
"Technology",
"Engineering"
] | 171 | [
"Spectroscopy stubs",
"Luminescence",
"Fluorescence",
"Spectrum (physical sciences)",
"Applied mathematics",
"Astronomy stubs",
"Measuring instruments",
"Computational science",
"Microscopes",
"Microscopy",
"Molecular physics stubs",
"Spectroscopy",
"Physical chemistry stubs"
] |
4,240,032 | https://en.wikipedia.org/wiki/Solenoid%20voltmeter | A solenoid voltmeter is a specific type of voltmeter electricians use to test electrical power circuits. It uses a solenoid coil to attract a spring-loaded plunger; the movement of the plunger is calibrated in terms of approximate voltage. It is more rugged than a D'arsonval movement, but neither as sensitive nor as precise.
Wiggy is the registered trademark for a common solenoid voltmeter used in North America derived from a device patent assigned to the Wigginton Company, US patent number 1,538,906.
Operation
Rather than using a D'Arsonval movement or digital electronics, the solenoid voltmeter simply uses a spring-loaded solenoid carrying a pointer (it might also be described as a form of moving iron meter). Greater voltage creates more magnetism pulling the solenoid's core in further against the spring loading, moving the pointer. A short scale converts the pointer's movement into the voltage reading. Solenoid voltmeters usually have a scale on each side of the pointer; one is calibrated for alternating current and one is calibrated for direct current. Only one "range" is provided and it usually extends from zero to about 600 volts.
A small permanent magnet rotor is usually mounted at the top of the meter. For DC, this magnet flips one way or the other, indicating by the exposed color (red or black) which lead is connected to positive. For AC, the rotor simply vibrates, indicating that the meter is connected to an AC circuit. Another form of tester uses a miniature neon lamp; the negative electrode glows, indicating polarity on DC circuits, or both electrodes glow, indicating AC.
Models made by some manufacturers include continuity test lights, which are energized by a battery within the tester. This is particularly advantageous when testing, for example, fuses in live circuits, since no switching is required to change from continuity mode to voltage detecting mode.
Comparison with moving coil meters
Solenoid voltmeters are extremely rugged and not very susceptible to damage through either rough handling or electrical overload, compared with more delicate but more precise instruments of the moving-coil D'arsonval type
For "go/no go" testing, there is no need to read the scale as application of AC power creates a perceivable vibration and sound within the meter. This feature makes the tester very handy in noisy, poorly illuminated, or very bright surroundings. The meter can be felt, the more it jumps the higher the voltage.
Solenoid voltmeters draw appreciable current in operation. When testing power supply circuits, a high-impedance connection (that is, a nearly open-circuit fault such as a burned switch contact or wire joint) in the power path might still allow enough voltage/current through to register on a high-impedance digital voltmeter, but it probably can't actuate the solenoid voltmeter. For use with high impedance circuit applications, however, they are not so good, as they draw appreciable current and therefore alter the voltage being measured. They can be used to test residual-current devices (GFCIs) because the current drawn trips most RCDs when the solenoid voltmeter is connected between the live and earth conductors.
Some manufacturers include a continuity test lamp function in a solenoid meter; these use the same probes as the voltage test function. This feature is useful when testing the status of contacts in energized circuits. The continuity light displays if the contact is closed, and the solenoid voltmeter shows voltage presence if open (and energized).
In contrast to multimeters, solenoid voltmeters have no other built-in functions (such as the ability to act as an ammeter, ohmmeter, or capacitance meter); they are just simple, easy-to-use power voltmeters. Solenoid voltmeters are useless on low-voltage circuits (for example, 12 volt circuits). The basic range of the voltmeter starts at around 90V (AC or DC).
Solenoid voltmeters are not precise. For example, there would be no reliably perceptible difference in the reading between 220 VAC and 240 VAC.
They are meant for intermittent operation. They draw a moderate amount of power from the circuit under test and can overheat if used for continuous monitoring.
The low impedance and low sensitivity of the tester may not show high-impedance connections to a voltage source, which can still source enough current to cause a shock hazard.
See also
Test light
Continuity tester
References
External links
All About Wiggy
Voltmeters
Electrical test equipment | Solenoid voltmeter | [
"Physics",
"Technology",
"Engineering"
] | 1,000 | [
"Voltmeters",
"Physical quantities",
"Electrical test equipment",
"Measuring instruments",
"Voltage"
] |
4,240,854 | https://en.wikipedia.org/wiki/Bipolaron | In physics, a bipolaron is a type of quasiparticle consisting of two polarons. In organic chemistry, it is a molecule or a part of a macromolecular chain containing two positive charges in a conjugated system.
Bipolarons in physics
In physics, a bipolaron is a bound pair of two polarons. An electron in a material may cause a distortion in the underlying lattice. The combination of electron and distortion (which may also be understood as a cloud of phonons) is known as a polaron (in part because the interaction between electron and lattice is via a polarization). When two polarons are close together, they can lower their energy by sharing the same distortions, which leads to an effective attraction between the polarons. If the interaction is sufficiently large, then that attraction leads to a bound bipolaron. For strong attraction, bipolarons may be small. Small bipolarons have integer spin and thus share some of the properties of bosons. If many bipolarons form without coming too close, they might be able to form a Bose–Einstein condensate. This has led to a suggestion that bipolarons could be a possible mechanism for high-temperature superconductivity. For example, they can lead to a very direct interpretation of the isotope effect.
Recently, bipolarons were predicted theorethically in a Bose-Einstein condensate. Two polarons interchange sound waves and they attract to each other, forming a bound-state when the strength coupling between the single polarons and the condensate is strong in comparison with the interactions of the host gas.
Bipolarons in organic chemistry
In organic chemistry, a bipolaron is a molecule or part of a macromolecular chain containing two positive charges in a conjugated system. The charges can be located in the centre of the chain or at its termini. Bipolarons and polarons are encountered in doped conducting polymers such as polythiophene.
It is possible to synthesize and isolate bipolaron model compounds for X-ray diffraction studies. The diamagnetic bis(triaryl)amine dication 2 in scheme 1 is prepared from the neutral precursor 1 in dichloromethane by reaction with 4 equivalents of antimony pentachloride. Two resonance structures exist for the dication. Structure 2a is a (singlet) diradical and 2b is the closed shell quinoid. The experimental bond lengths for the central vinylidene group in 2 are 141 pm and 137 pm compared to 144 pm and 134 pm for the precursor 1 implying some contribution from the quinoid structure.
On the other hand, when a thiophene unit is added to the core in the structure depicted in scheme 2, these bond lengths are identical (around 138 pm) making it a true hybrid.
See also
Quinonoid zwitterions
References
Ions
Quasiparticles | Bipolaron | [
"Physics",
"Materials_science"
] | 598 | [
"Quasiparticles",
"Subatomic particles",
"Condensed matter physics",
"Matter"
] |
4,242,000 | https://en.wikipedia.org/wiki/Power%20system%20protection | Power system protection is a branch of electrical power engineering that deals with the protection of electrical power systems from faults through the disconnection of faulted parts from the rest of the electrical network. The objective of a protection scheme is to keep the power system stable by isolating only the components that are under fault, whilst leaving as much of the network as possible in operation. The devices that are used to protect the power systems from faults are called protection devices.
Components
Protection systems usually comprise five components
Current and voltage transformers to step down the high voltages and currents of the electrical power system to convenient levels for the relays to deal with
Protective relays to sense the fault and initiate a trip, or disconnection, order
Circuit breakers or RCDs to open/close the system based on relay and autorecloser commands
Batteries to provide power in case of power disconnection in the system
Communication channels to allow analysis of current and voltage at remote terminals of a line and to allow remote tripping of equipment.
For parts of a distribution system, fuses are capable of both sensing and disconnecting faults.
Failures may occur in each part, such as insulation failure, fallen or broken transmission lines, incorrect operation of circuit breakers, short circuits and open circuits. Protection devices are installed with the aims of protection of assets and ensuring continued supply of energy.
Switchgear is a combination of electrical disconnect switches, fuses or circuit breakers used to control, protect and isolate electrical equipment. Switches are safe to open under normal load current (some switches are not safe to operate under normal or abnormal conditions), while protective devices are safe to open under fault current. Very important equipment may have completely redundant and independent protective systems, while a minor branch distribution line may have very simple low-cost protection.
Types of protection
High-voltage transmission network
Protection of the transmission and distribution system serves two functions: protection of the plant and protection of the public (including employees). At a basic level, protection disconnects equipment that experiences an overload or a short to earth. Some items in substations such as transformers might require additional protection based on temperature or gas pressure, among others.
Generator sets
In a power plant, the protective relays are intended to prevent damage to alternators or to the transformers in case of abnormal conditions of operation, due to internal failures, as well as insulating failures or regulation malfunctions. Such failures are unusual, so the protective relays have to operate very rarely. If a protective relay fails to detect a fault, the resulting damage to the alternator or to the transformer might require costly equipment repairs or replacement, as well as income loss from the inability to produce and sell energy.
Overload and back-up for distance (overcurrent)
Overload protection requires a current transformer which simply measures the current in a circuit and compares it to the predetermined value. There are two types of overload protection: instantaneous overcurrent (IOC) and time overcurrent (TOC). Instantaneous overcurrent requires that the current exceeds a predetermined level for the circuit breaker to operate. Time overcurrent protection operates based on a current vs time curve. Based on this curve, if the measured current exceeds a given level for the preset amount of time, the circuit breaker or fuse will operate. The function of both types is explained in .
Earth fault/ground fault
Earth fault protection also requires current transformers and senses an imbalance in a three-phase circuit. Normally the three phase currents are in balance, i.e. roughly equal in magnitude. If one or two phases become connected to earth via a low impedance path, their magnitudes will increase dramatically, as will current imbalance. If this imbalance exceeds a pre-determined value, a circuit breaker should operate. Restricted earth fault protection is a type of earth fault protection which looks for earth fault between two sets of current transformers (hence restricted to that zone).
Distance (impedance relay)
Distance protection detects both voltage and current. A fault on a circuit will generally create a sag in the voltage level. If the ratio of voltage to current measured at the relay terminals, which equates to an impedance, lands within a predetermined level the circuit breaker will operate. This is useful for reasonably long lines, lines longer than 10 miles, because their operating characteristics are based on the line characteristics. This means that when a fault appears on the line the impedance setting in the relay is compared to the apparent impedance of the line from the relay terminals to the fault. If the relay setting is determined to be below the apparent impedance it is determined that the fault is within the zone of protection. When the transmission line length is too short, less than 10 miles, distance protection becomes more difficult to coordinate. In these instances the best choice of protection is current differential protection.
Back-up
The objective of protection is to remove only the affected portion of plant and nothing else. A circuit breaker or protection relay may fail to operate. In important systems, a failure of primary protection will usually result in the operation of back-up protection. Remote back-up protection will generally remove both the affected and unaffected items of plant to clear the fault. Local back-up protection will remove the affected items of the plant to clear the fault.
Low-voltage networks
The low-voltage network generally relies upon fuses or low-voltage circuit breakers to remove both overload and earth faults.
Cybersecurity
The bulk system which is a large interconnected electrical system including transmission and control system is experiencing new cybersecurity threats every day. (“Electric Grid Cybersecurity,” 2019). Most of these attacks are aiming the control systems in the grids. These control systems are connected to the internet and makes it easier for hackers to attack them. These attacks can cause damage to equipment and limit the utility professionals ability to control the system.
Coordination
Protective device coordination is the process of determining the "best fit" timing of current interruption when abnormal electrical conditions occur. The goal is to minimize an outage to the greatest extent possible. Historically, protective device coordination was done on translucent log–log paper. Modern methods normally include detailed computer based analysis and reporting.
Protection coordination is also handled through dividing the power system into protective zones. If a fault were to occur in a given zone, necessary actions will be executed to isolate that zone from the entire system. Zone definitions account for generators, buses, transformers, transmission and distribution lines, and motors. Additionally, zones possess the following features: zones overlap, overlap regions denote circuit breakers, and all circuit breakers in a given zone with a fault will open in order to isolate the fault. Overlapped regions are created by two sets of instrument transformers and relays for each circuit breaker. They are designed for redundancy to eliminate unprotected areas; however, overlapped regions are devised to remain as small as possible such that when a fault occurs in an overlap region and the two zones which encompass the fault are isolated, the sector of the power system which is lost from service is still small despite two zones being isolated.
Disturbance-monitoring equipment
Disturbance-monitoring equipment (DME) monitors and records system data pertaining to a fault. DME accomplish three main purposes:
model validation,
disturbance investigation, and
assessment of system protection performance.
DME devices include:
Sequence of event recorders, which record equipment response to the event
Fault recorders, which record actual waveform data of the system primary voltages and currents
Dynamic disturbance recorders (DDRs), which record incidents that portray power system behavior during dynamic events such as low frequency (0.1 Hz – 3 Hz) oscillations and abnormal frequency or voltage excursions
Performance measures
Protection engineers define dependability as the tendency of the protection system to operate correctly for in-zone faults. They define security as the tendency not to operate for out-of-zone faults. Both dependability and security are reliability issues. Fault tree analysis is one tool with which a protection engineer can compare the relative reliability of proposed protection schemes. Quantifying protection reliability is important for making the best decisions on improving a protection system, managing dependability versus security tradeoffs, and getting the best results for the least money. A quantitative understanding is essential in the competitive utility industry.
Reliability: Devices must function consistently when fault conditions occur, regardless of possibly being idle for months or years. Without this reliability, systems may cause costly damages.
Selectivity: Devices must avoid unwarranted, false trips.
Speed: Devices must function quickly to reduce equipment damage and fault duration, with only very precise intentional time delays.
Sensitivity: Devices must detect even the smallest value of faults and respond.
Economy: Devices must provide maximum protection at minimum cost.
Simplicity: Devices must minimize protection circuitry and equipment.
Reliability: Dependability vs Security
There are two aspects of reliable operation of protection systems: dependability and security. Dependability is the ability of the protection system to operate when called upon to remove a faulted element from the power system. Security is the ability of the protection system to restrain itself from operating during an external fault. Choosing the appropriate balance between security and dependability in designing the protection system requires engineering judgement and varies on a case-by-case basis.
See also
Fault current limiter
Network analyzer (AC power)
Prospective short-circuit current
ANSI device numbers
Notes
References
http://perso.numericable.fr/michlami protection and monitoring of the electrical energy transmission network
Over-current protection devices
Power engineering | Power system protection | [
"Engineering"
] | 1,954 | [
"Power engineering",
"Electrical engineering",
"Energy engineering"
] |
4,242,960 | https://en.wikipedia.org/wiki/National%20Fluid%20Power%20Association | The National Fluid Power Association (NFPA) is an American 501(c)6 industry trade association, founded in 1953.
The NFPA's mission is to serve as a forum where all fluid power channel partners work together to advance fluid power technology, strengthen the fluid power industry, and foster members' success. NFPA members include more than 315 manufacturers of fluid power systems and components, fluid power distributors, suppliers to the fluid power industry, educators and researchers.
References
External links
National Fluid Power Association
Fluid Power Distributors Association
Fluid Power Education Foundation
Fluid Power Society
Hydraulic engineering organizations
Trade associations based in the United States | National Fluid Power Association | [
"Engineering"
] | 126 | [
"Hydraulic engineering organizations",
"Civil engineering organizations"
] |
8,877,643 | https://en.wikipedia.org/wiki/Misner%20space | Misner space is an abstract mathematical spacetime, first described by Charles W. Misner. It is also known as the Lorentzian orbifold . It is a simplified, two-dimensional version of the Taub–NUT spacetime. It contains a non-curvature singularity and is an important counterexample to various hypotheses in general relativity.
Michio Kaku develops the following analogy for understanding the concept: "Misner space is an idealized space in which a room, for example, becomes the entire universe. For example, every point on the left wall of the room is identical to the corresponding point on the right wall, such that if you were to walk toward the left wall you will walk through the wall and appear from the right wall. This suggests that the left and right wall are joined, in some sense, as in a cylinder. The opposite walls are thus all identified with each other, and the ceiling is likewise identified with the floor. Misner space is often studied because it has the same topology as a wormhole but is much simpler to handle mathematically. If the walls move, then time travel might be possible within the Misner universe."
Metric
The simplest description of Misner space is to consider two-dimensional Minkowski space with the metric
with the identification of every pair of spacetime points by a constant boost
It can also be defined directly on the cylinder manifold with coordinates by the metric
The two coordinates are related by the map
and
Causality
Misner space is a standard example for the study of causality since it contains both closed timelike curves and a compactly generated Cauchy horizon, while still being flat (since it is just Minkowski space). With the coordinates , the loop defined by , with tangent vector , has the norm , making it a closed null curve. This is the chronology horizon : there are no closed timelike curves in the region , while every point admits a closed timelike curve through it in the region .
This is due to the tipping of the light cones which, for , remains above lines of constant but will open beyond that line for , causing any loop of constant to be a closed timelike curve.
Chronology protection
Misner space was the first spacetime where the notion of chronology protection was used for quantum fields, by showing that in the semiclassical approximation, the expectation value of the stress-energy tensor for the vacuum is divergent.
References
Further reading
General relativity | Misner space | [
"Physics"
] | 499 | [
"General relativity",
"Theory of relativity"
] |
8,884,972 | https://en.wikipedia.org/wiki/Kirchhoff%20equations | In fluid dynamics, the Kirchhoff equations, named after Gustav Kirchhoff, describe the motion of a rigid body in an ideal fluid.
where and are the angular and linear velocity vectors at the point , respectively; is the moment of inertia tensor, is the body's mass; is
a unit normal vector to the surface of the body at the point ;
is a pressure at this point; and are the hydrodynamic
torque and force acting on the body, respectively;
and likewise denote all other torques and forces acting on the
body. The integration is performed over the fluid-exposed portion of the
body's surface.
If the body is completely submerged body in an infinitely large volume of irrotational, incompressible, inviscid fluid, that is at rest at infinity, then the vectors and can be found via explicit integration, and the dynamics of the body is described by the Kirchhoff – Clebsch equations:
Their first integrals read
Further integration produces explicit expressions for position and velocities.
References
Kirchhoff G. R. Vorlesungen ueber Mathematische Physik, Mechanik. Lecture 19. Leipzig: Teubner. 1877.
Lamb, H., Hydrodynamics. Sixth Edition Cambridge (UK): Cambridge University Press. 1932.
Mechanics
Classical mechanics
Rigid bodies
Gustav Kirchhoff | Kirchhoff equations | [
"Physics",
"Chemistry",
"Engineering"
] | 285 | [
"Classical mechanics",
"Mechanics",
"Mechanical engineering",
"Fluid dynamics stubs",
"Fluid dynamics"
] |
5,650,682 | https://en.wikipedia.org/wiki/The%20Regenerative%20Medicine%20Institute | The Regenerative Medicine Institute (REMEDI), was established in 2003 as a Centre for Science, Technology & Engineering in collaboration with National University of Ireland, Galway. It obtained an award of €14.9 million from Science Foundation Ireland over five years.
It conducts basic research and applied research in regenerative medicine, an emerging field that combines the technologies of gene therapy and adult stem cell therapy. The goal is to use cells and genes to regenerate healthy tissues that can be used to repair or replace other tissues and organs in a minimally invasive approach.
Centres for Science, Engineering & Technology help link scientists and engineers in partnerships across academia and industry to address crucial research questions, foster the development of new and existing Irish-based technology companies, attract industry that could make an important contribution to Ireland and its economy, and expand educational and career opportunities in Ireland in science and engineering. CSETs must exhibit outstanding research quality, intellectual breadth, active collaboration, flexibility in responding to new research opportunities, and integration of research and education in the fields that SFI supports.
References
External links
Regenerative Medicine Institute (REMEDI)
Science Foundation Ireland
National University of Ireland, Galway
Medical research institutes in the Republic of Ireland
Biotechnology organizations
Bioethics research organizations
2003 establishments in Ireland
Scientific organizations established in 2003 | The Regenerative Medicine Institute | [
"Engineering",
"Biology"
] | 263 | [
"Biotechnology organizations"
] |
5,654,027 | https://en.wikipedia.org/wiki/Synapse | In the nervous system, a synapse is a structure that allows a neuron (or nerve cell) to pass an electrical or chemical signal to another neuron or a target effector cell. Synapses can be classified as either chemical or electrical, depending on the mechanism of signal transmission between neurons. In the case of electrical synapses, neurons are coupled bidirectionally with each other through gap junctions and have a connected cytoplasmic milieu. These types of synapses are known to produce synchronous network activity in the brain, but can also result in complicated, chaotic network level dynamics. Therefore, signal directionality cannot always be defined across electrical synapses.
Synapses are essential for the transmission of neuronal impulses from one neuron to the next, playing a key role in enabling rapid and direct communication by creating circuits. In addition, a synapse serves as a junction where both the transmission and processing of information occur, making it a vital means of communication between neurons.
At the synapse, the plasma membrane of the signal-passing neuron (the presynaptic neuron) comes into close apposition with the membrane of the target (postsynaptic) cell. Both the presynaptic and postsynaptic sites contain extensive arrays of molecular machinery that link the two membranes together and carry out the signaling process. In many synapses, the presynaptic part is located on the terminals of axons and the postsynaptic part is located on a dendrite or soma. Astrocytes also exchange information with the synaptic neurons, responding to synaptic activity and, in turn, regulating neurotransmission. Synapses (at least chemical synapses) are stabilized in position by synaptic adhesion molecules (SAMs) projecting from both the pre- and post-synaptic neuron and sticking together where they overlap; SAMs may also assist in the generation and functioning of synapses. Moreover, SAMs coordinate the formation of synapses, with various types working together to achieve the remarkable specificity of synapses. In essence, SAMs function in both excitatory and inhibitory synapses, likely serving as the mediator for signal transmission.
History
Santiago Ramón y Cajal proposed that neurons are not continuous throughout the body, yet still communicate with each other, an idea known as the neuron doctrine. The word "synapse" was introduced in 1897 by the English neurophysiologist Charles Sherrington in Michael Foster's Textbook of Physiology. Sherrington struggled to find a good term that emphasized a union between two separate elements, and the actual term "synapse" was suggested by the English classical scholar Arthur Woollgar Verrall, a friend of Foster. The word was derived from the Greek synapsis (), meaning "conjunction", which in turn derives from synaptein (), from syn () "together" and haptein () "to fasten".
However, while the synaptic gap remained a theoretical construct, and was sometimes reported as a discontinuity between contiguous axonal terminations and dendrites or cell bodies, histological methods using the best light microscopes of the day could not visually resolve their separation which is now known to be about 20 nm. It needed the electron microscope in the 1950s to show the finer structure of the synapse with its separate, parallel pre- and postsynaptic membranes and processes, and the cleft between the two.
Types
Chemical and electrical synapses are two ways of synaptic transmission.
In a chemical synapse, electrical activity in the presynaptic neuron is converted (via the activation of voltage-gated calcium channels) into the release of a chemical called a neurotransmitter that binds to receptors located in the plasma membrane of the postsynaptic cell. The neurotransmitter may initiate an electrical response or a secondary messenger pathway that may either excite or inhibit the postsynaptic neuron. Chemical synapses can be classified according to the neurotransmitter released: glutamatergic (often excitatory), GABAergic (often inhibitory), cholinergic (e.g. vertebrate neuromuscular junction), and adrenergic (releasing norepinephrine). Because of the complexity of receptor signal transduction, chemical synapses can have complex effects on the postsynaptic cell.
In an electrical synapse, the presynaptic and postsynaptic cell membranes are connected by special channels called gap junctions that are capable of passing an electric current, causing voltage changes in the presynaptic cell to induce voltage changes in the postsynaptic cell. In fact, gap junctions facilitate the direct flow of electrical current without the need for neurotransmitters, as well as small molecules like calcium. Thus, the main advantage of an electrical synapse is the rapid transfer of signals from one cell to the next.
Mixed chemical electrical synapses are synaptic sites that feature both a gap junction and neurotransmitter release. This combination allows a signal to have both a fast component (electrical) and a slow component (chemical).
The formation of neural circuits in nervous systems appears to heavily depend on the crucial interactions between chemical and electrical synapses. Thus these interactions govern the generation of synaptic transmission. Synaptic communication is distinct from an ephaptic coupling, in which communication between neurons occurs via indirect electric fields. An autapse is a chemical or electrical synapse that forms when the axon of one neuron synapses onto dendrites of the same neuron.
Excitatory and inhibitory
Excitatory synapse: Enhances the probability of depolarization in postsynaptic neurons and the initiation of an action potential.
Inhibitory synapse: Diminishes the probability of depolarization in postsynaptic neurons and the initiation of an action potential.
An influx of Na+ driven by excitatory neurotransmitters opens cation channels, depolarizing the postsynaptic membrane toward the action potential threshold. In contrast, inhibitory neurotransmitters cause the postsynaptic membrane to become less depolarized by opening either Cl- or K+ channels, reducing firing. Depending on their release location, the receptors they bind to, and the ionic circumstances they encounter, various transmitters can be either excitatory or inhibitory. For instance, acetylcholine can either excite or inhibit depending on the type of receptors it binds to. For example, glutamate serves as an excitatory neurotransmitter, in contrast to GABA, which acts as an inhibitory neurotransmitter. Additionally, dopamine is a neurotransmitter that exerts dual effects, displaying both excitatory and inhibitory impacts through binding to distinct receptors.
The membrane potential prevents Cl- from entering the cell, even when its concentration is much higher outside than inside. The reversal potential for Cl- in many neurons is quite negative, nearly equal to the resting potential. Opening Cl- channels tends to buffer the membrane potential, but this effect is countered when the membrane starts to depolarize, allowing more negatively charged Cl- ions to enter the cell. Consequently, it becomes more difficult to depolarize the membrane and excite the cell when Cl- channels are open. Similar effects result from the opening of K+ channels. The significance of inhibitory neurotransmitters is evident from the effects of toxins that impede their activity. For instance, strychnine binds to glycine receptors, blocking the action of glycine and leading to muscle spasms, convulsions, and death.
Interfaces
Synapses can be classified by the type of cellular structures serving as the pre- and post-synaptic components. The vast majority of synapses in the mammalian nervous system are classical axo-dendritic synapses (axon synapsing upon a dendrite), however, a variety of other arrangements exist. These include but are not limited to axo-axonic, dendro-dendritic, axo-secretory, axo-ciliary, somato-dendritic, dendro-somatic, and somato-somatic synapses.
In fact, the axon can synapse onto a dendrite, onto a cell body, or onto another axon or axon terminal, as well as into the bloodstream or diffusely into the adjacent nervous tissue.
Conversion of chemical into electrical signals
Neurotransmitters are tiny signal molecules stored in membrane-enclosed synaptic vesicles and released via exocytosis. Indeed, a change in electrical potential in the presynaptic cell triggers the release of these molecules. By attaching to transmitter-gated ion channels, the neurotransmitter causes an electrical alteration in the postsynaptic cell and rapidly diffuses across the synaptic cleft. Once released, the neurotransmitter is swiftly eliminated, either by being absorbed by the nerve terminal that produced it, taken up by nearby glial cells, or broken down by specific enzymes in the synaptic cleft. Numerous Na+-dependent neurotransmitter carrier proteins recycle the neurotransmitters and enable the cells to maintain rapid rates of release.
At chemical synapses, transmitter-gated ion channels play a vital role in rapidly converting extracellular chemical impulses into electrical signals. These channels are located in the postsynaptic cell's plasma membrane at the synapse region, and they temporarily open in response to neurotransmitter molecule binding, causing a momentary alteration in the membrane's permeability. Additionally, transmitter-gated channels are comparatively less sensitive to the membrane potential than voltage-gated channels, which is why they are unable to generate self-amplifying excitement on their own. However, they result in graded variations in membrane potential due to local permeability, influenced by the amount and duration of neurotransmitter released at the synapse.
Recently, mechanical tension, a phenomenon never thought relevant to synapse function has been found to be required for those on hippocampal neurons to fire.
Release of neurotransmitters
Neurotransmitters bind to ionotropic receptors on postsynaptic neurons, either causing their opening or closing. The variations in the quantities of neurotransmitters released from the presynaptic neuron may play a role in regulating the effectiveness of synaptic transmission. In fact, the concentration of cytoplasmic calcium is involved in regulating the release of neurotransmitters from presynaptic neurons.
The chemical transmission involves several sequential processes:
Synthesizing neurotransmitters within the presynaptic neuron.
Loading the neurotransmitters into secretory vesicles.
Controlling the release of neurotransmitters into the synaptic cleft.
Binding of neurotransmitters to postsynaptic receptors.
Ceasing the activity of the released neurotransmitters.
Synaptic polarization
The function of neurons depends upon cell polarity. The distinctive structure of nerve cells allows action potentials to travel directionally (from dendrites to cell body down the axon), and for these signals to then be received and carried on by post-synaptic neurons or received by effector cells. Nerve cells have long been used as models for cellular polarization, and of particular interest are the mechanisms underlying the polarized localization of synaptic molecules. PIP2 signaling regulated by IMPase plays an integral role in synaptic polarity.
Phosphoinositides (PIP, PIP2, and PIP3) are molecules that have been shown to affect neuronal polarity. A gene (ttx-7) was identified in Caenorhabditis elegans that encodes myo-inositol monophosphatase (IMPase), an enzyme that produces inositol by dephosphorylating inositol phosphate. Organisms with mutant ttx-7 genes demonstrated behavioral and localization defects, which were rescued by expression of IMPase. This led to the conclusion that IMPase is required for the correct localization of synaptic protein components. The egl-8 gene encodes a homolog of phospholipase Cβ (PLCβ), an enzyme that cleaves PIP2. When ttx-7 mutants also had a mutant egl-8 gene, the defects caused by the faulty ttx-7 gene were largely reversed. These results suggest that PIP2 signaling establishes polarized localization of synaptic components in living neurons.
Presynaptic modulation
Modulation of neurotransmitter release by G-protein-coupled receptors (GPCRs) is a prominent presynaptic mechanism for regulation of synaptic transmission. The activation of GPCRs located at the presynaptic terminal, can decrease the probability of neurotransmitter release. This presynaptic depression involves activation of Gi/o-type G-proteins that mediate different inhibitory mechanisms, including inhibition of voltage-gated calcium channels, activation of potassium channels, and direct inhibition of the vesicle fusion process.
Endocannabinoids, synthesized in and released from postsynaptic neuronal elements and their cognate receptors, including the (GPCR) CB1 receptor located at the presynaptic terminal, are involved in this modulation by a retrograde signaling process, in which these compounds are synthesized in and released from postsynaptic neuronal elements and travel back to the presynaptic terminal to act on the CB1 receptor for short-term or long-term synaptic depression, that causes a short or long lasting decrease in neurotransmitter release.
Effects of drugs on ligand-gated ion channels
Drugs have long been considered crucial targets for transmitter-gated ion channels. The majority of medications utilized to treat schizophrenia, anxiety, depression, and sleeplessness work at chemical synapses, and many of these pharmaceuticals function by binding to transmitter-gated channels. For instance, some drugs like barbiturates and tranquilizers bind to GABA receptors and enhance the inhibitory effect of GABA neurotransmitter. Thus, reduced concentration of GABA enables the opening of Cl- channels.
Furthermore, psychoactive drugs could potentially target many other synaptic signalling machinery components. In fact, numerous neurotransmitters are released by Na+-driven carriers and are subsequently removed from the synaptic cleft. By inhibiting such carriers, synaptic transmission is strengthened as the action of the transmitter is prolonged. For example, Prozac is an antidepressant medication that works by preventing the absorption of serotonin neurotransmitter. Also, other antidepressants operate by inhibiting the reabsorption of both serotonin and norepinephrine.
Biogenesis
In nerve terminals, synaptic vesicles are produced quickly to compensate for their rapid depletion during neurotransmitter release. Their biogenesis involves segregating synaptic vesicle membrane proteins from other cellular proteins and packaging those distinct proteins into vesicles of appropriate size. Besides, it entails the endocytosis of synaptic vesicle membrane proteins from the plasma membrane.
Synaptoblastic and synaptoclastic refer to synapse-producing and synapse-removing activities within the biochemical signalling chain. This terminology is associated with the Bredesen Protocol for treating Alzheimer's disease, which conceptualizes Alzheimer's as an imbalance between these processes. As of October 2023, studies concerning this protocol remain small and few results have been obtained within a standardized control framework.
Role in memory
Potentiation and depression
It is widely accepted that the synapse plays a key role in the formation of memory. The stability of long-term memory can persist for many years; nevertheless, synapses, the neurological basis of memory, are very dynamic. The formation of synaptic connections significantly depends on activity-dependent synaptic plasticity observed in various synaptic pathways. Indeed, the connection between memory formation and alterations in synaptic efficacy enables the reinforcement of neuronal interactions between neurons. As neurotransmitters activate receptors across the synaptic cleft, the connection between the two neurons is strengthened when both neurons are active at the same time, as a result of the receptor's signaling mechanisms. The strength of two connected neural pathways is thought to result in the storage of information, resulting in memory. This process of synaptic strengthening is known as long-term potentiation (LTP).
By altering the release of neurotransmitters, the plasticity of synapses can be controlled in the presynaptic cell. The postsynaptic cell can be regulated by altering the function and number of its receptors. Changes in postsynaptic signaling are most commonly associated with a N-methyl-d-aspartic acid receptor (NMDAR)-dependent LTP and long-term depression (LTD) due to the influx of calcium into the post-synaptic cell, which are the most analyzed forms of plasticity at excitatory synapses.
Mechanism of protein kinase
Moreover, Ca2+/calmodulin (CaM)-dependent protein kinase II (CaMKII) is best recognized for its roles in the brain, particularly in the neocortex and hippocampal regions because it serves as a ubiquitous mediator of cellular Ca2+ signals. CaMKII is abundant in the nervous system, mainly concentrated in the synapses in the nerve cells. Indeed, CaMKII has been definitively identified as a key regulator of cognitive processes, such as learning, and neural plasticity. The first concrete experimental evidence for the long-assumed function of CaMKII in memory storage was demonstrated
While Ca2+/CaM binding stimulates CaMKII activity, Ca2+-independent autonomous CaMKII activity can also be produced by a number of other processes. CaMKII becomes active by autophosphorylating itself upon Ca2+/calmodulin binding. CaMKII is still active and phosphorylates itself even after Ca2+ is cleaved; as a result, the brain stores long-term memories using this mechanism. Nevertheless, when the CaMKII enzyme is dephosphorylated by a phosphatase enzyme, it becomes inactive, and memories are lost. Hence, CaMKII plays a vital role in both the induction and maintenance of LTP.
Experimental models
For technical reasons, synaptic structure and function have been historically studied at unusually large model synapses, for example:
Squid giant synapse
Neuromuscular junction (NMJ), a cholinergic synapse in vertebrates, glutamatergic in insects
Ciliary calyx in the ciliary ganglion of chicks
Calyx of Held in the brainstem
Ribbon synapse in the retina
Schaffer collateral synapses in the hippocampus. These synapses are small, but their pre- and postsynaptic neurons are well separated (CA3 and CA1, respectively).
Synapses and diseases
Synapses function as ensembles within particular brain networks to control the amount of neuronal activity, which is essential for memory, learning, and behavior. Consequently, synaptic disruptions might have negative effects. In fact, alterations in cell-intrinsic molecular systems or modifications to environmental biochemical processes can lead to synaptic dysfunction. The synapse is the primary unit of information transfer in the nervous system, and correct synaptic contact creation during development is essential for normal brain function. In addition, several mutations have been connected to neurodevelopmental disorders, and that compromised function at different synapse locations is a hallmark of neurodegenerative diseases.
Synaptic defects are causally associated with early appearing neurological diseases, including autism spectrum disorders (ASD), schizophrenia (SCZ), and bipolar disorder (BP). On the other hand, in late-onset degenerative pathologies, such as Alzheimer's (AD), Parkinson's (PD), and Huntington's (HD) diseases, synaptopathy is thought to be the inevitable end-result of an ongoing pathophysiological cascade. These diseases are identified by a gradual loss in cognitive and behavioral function and a steady loss of brain tissue. Moreover, these deteriorations have been mostly linked to the gradual build-up of protein aggregates in neurons, the composition of which may vary based on the pathology; all have the same deleterious effects on neuronal integrity. Furthermore, the high number of mutations linked to synaptic structure and function, as well as dendritic spine alterations in post-mortem tissue, has led to the association between synaptic defects and neurodevelopmental disorders, such as ASD and SCZ, characterized by abnormal behavioral or cognitive phenotypes.
Nevertheless, due to limited access to human tissue at late stages and a lack of thorough assessment of the essential components of human diseases in the available experimental animal models, it has been difficult to fully grasp the origin and role of synaptic dysfunction in neurological disorders.
Additional images
See also
Active zone
Autapse
Cooperative synapse formation
Exocytosis
Immunological synapse
Neurotransmitter vesicle
Neurexin
Postsynaptic density
Synaptopathy
References
Signal transduction | Synapse | [
"Chemistry",
"Biology"
] | 4,570 | [
"Biochemistry",
"Neurochemistry",
"Signal transduction"
] |
5,654,154 | https://en.wikipedia.org/wiki/Immunological%20synapse | In immunology, an immunological synapse (or immune synapse) is the interface between an antigen-presenting cell or target cell and a lymphocyte such as a T cell, B cell, or natural killer cell. The interface was originally named after the neuronal synapse, with which it shares the main structural pattern. An immunological synapse consists of molecules involved in T cell activation, which compose typical patterns—activation clusters. Immunological synapses are the subject of much ongoing research.
Structure and function
The immune synapse is also known as the supramolecular activation cluster or SMAC. This structure is composed of concentric rings each containing segregated clusters of proteins—often referred to as the bull’s-eye model of the immunological synapse:
c-SMAC (central-SMAC) composed of the θ isoform of protein kinase C, CD2, CD4, CD8, CD28, Lck, and Fyn.
p-SMAC (peripheral-SMAC) within which the lymphocyte function-associated antigen-1 (LFA-1) and the cytoskeletal protein talin are clustered.
d-SMAC (distal-SMAC) enriched in CD43 and CD45 molecules.
New investigations, however, have shown that a "bull’s eye" is not present in all immunological synapses. For example, different patterns appear in the synapse between a T-cell and a dendritic cell.
This complex as a whole is postulated to have several functions including but not limited to:
Regulation of lymphocyte activation
Transfer of peptide-MHC complexes from APCs to lymphocytes
Directing secretion of cytokines or lytic granules
Recent research has proposed a striking parallel between the immunological synapse and the primary cilium based mainly on similar actin rearrangement, orientation of the centrosome towards the structure and involvement of similar transport molecules (such as IFT20, Rab8, Rab11). This structural and functional homology is the topic of ongoing research.
Formation
The initial interaction occurs between LFA-1 present in the p-SMAC of a T-cell, and non-specific adhesion molecules (such as ICAM-1 or ICAM-2) on a target cell. When bound to a target cell, the T-cell can extend pseudopodia and scan the surface of target cell to find a specific peptide:MHC complex.
The process of formation begins when the T-cell receptor (TCR) binds to the peptide:MHC complex on the antigen-presenting cell and initiates signaling activation through formation of microclusters/lipid rafts. Specific signaling pathways lead to polarization of the T-cell by orienting its centrosome toward the site of the immunological synapse. The symmetric centripetal actin flow is the basis of formation of the p-SNAP ring. The accumulation and polarization of actin is triggered by TCR/CD3 interactions with integrins and small GTPases (such as Rac1 or Cdc42). These interactions activate large multi-molecular complexes (containing WAVE (Scar), HSP300, ABL2, SRA1, and NAP1 and others) to associate with Arp2/3, which directly promotes actin polymerization. As actin is accumulated and reorganized, it promotes clustering of TCRs and integrins. The process thereby upregulates itself via positive feedback.
Some parts of this process may differ in CD4+ and CD8+ cells. For example, synapse formation is quick in CD8+ T cells, because for CD8+ T cells it is fundamental to eliminate the pathogen quickly. In CD4+ T cells, however, the whole process of the immunological synapse formation can take up to 6 hours.
In CD8+ T cells, the synapse formation leads to killing of the target cell via secretion of cytolytic enzymes. CD8+ T lymphocytes contain lytic granules – specialized secretory lysosomes filled with perforin, granzymes, lysosomal hydrolases (for example cathepsins B and D, β-hexosaminidase) and other cytolytic effector proteins. Once these proteins are delivered to the target cell, they induce its apoptosis. The effectivity of killing of the target cell depends on the strength of the TCR signal. Even after receiving weak or short-lived signals, the MTOC polarizes towards the immunological synapse, but in that case the lytic granules are not trafficked and therefore the killing effect is missing or poor.
NK-cell synapse
NK cells are known to form synapses with cytolytic effect towards the target cell. In the initiation step, NK cell approaches the target cell, either accidentally or intentionally due to the chemotactic signalling. Firstly, the sialyl Lewis X present on the surface of target cell is recognized by CD2 on NK cell. If the KIR receptors of NK cell find their cognate antigen on the surface of target cell, formation of the lytic synapse is inhibited. If such signal is missing, a tight adhesion via LFA1 and MAC1 is promoted and enhanced by additional signals such as CD226-ligand and CD96-CD155 interactions.
Lytic granules are secretory organelles filled with perforin, granzymes and other cytolytic enzymes. After initiation of the cell-cell contact, the lytic granules of NK cells move around the microtubules towards the centrosome, which also relocalizes towards the site of synapse. Then, the contents of lytic granules is released and via vesicles with SNARE proteins transferred to the target cell.
Inhibitory immunological synapse of NK cells
When an NK cell encounters a self cell, it forms a so-called inhibitory immunological synapse to prevent unwanted cytolysis of target cell. In this process, the killer-cell immunoglobulin-like receptors (KIRs) containing long cytoplasmic tails with immunoreceptor tyrosine-based inhibitory motifs (ITIMs) are clustered in the site of synapse, bind their ligand on the surface of target cell and form the supramolecular inhibitory cluster (SMIC). SMIC then acts to prevent rearrangement of actin, block the recruitment of activatory receptors to the site of synapse and finally, promote detachment from the target cell. This process is essential in protecting NK cells from killing self cells.
History
Immunological synapses were first discovered by Abraham Kupfer at the National Jewish Medical and Research Center in Denver. Their name was coined by Michael Dustin at NYU who studied them in further detail. Daniel M. Davis and Jack Strominger showed structured immune synapses for a different lymphocyte, the Natural Killer cell, and published this around the same time. Abraham Kupfer first presented his findings during a Keystone Symposia in 1995, when he showed three-dimensional images of immune cells interacting with one another. Key molecules in the synapse are the T cell receptor and its counterpart the major histocompatibility complex (MHC). Also important are LFA-1, ICAM-1, CD28, and CD80/CD86.
References
External links
Immunological Synapse - Cell Centered Database
Immune system | Immunological synapse | [
"Biology"
] | 1,606 | [
"Immune system",
"Organ systems"
] |
5,654,239 | https://en.wikipedia.org/wiki/Bone%20morphogenetic%20protein%203 | Bone morphogenetic protein 3, also known as osteogenin, is a protein in humans that is encoded by the BMP3 gene.
The protein encoded by this gene is a member of the transforming growth factor beta superfamily. It, unlike other bone morphogenetic proteins (BMP's) inhibits the ability of other BMP's to induce bone and cartilage development. It is a disulfide-linked homodimer. It negatively regulates bone density. BMP3 is an antagonist to other BMP's in the differentiation of osteogenic progenitors.
It is highly expressed in fractured tissues.
Cancer
BMP3 is hypermethylated in many cases of colorectal cancer (CRC) and hence along with other hypermethylated genes, may be used as a biomarker to detect early stage CRC.
References
External links
Further reading
Bone morphogenetic protein
Developmental genes and proteins
TGFβ domain | Bone morphogenetic protein 3 | [
"Biology"
] | 200 | [
"Induced stem cells",
"Developmental genes and proteins"
] |
5,654,541 | https://en.wikipedia.org/wiki/Software%20security%20assurance | Software security assurance is a process that helps design and implement software that protects the data and resources contained in and controlled by that software. Software is itself a resource and thus must be afforded appropriate security.
What is software security assurance?
Software Security Assurance (SSA) is the process of ensuring that software is designed to operate at a level of security that is consistent with the potential harm that could result from the loss, inaccuracy, alteration, unavailability, or misuse of the data and resources that it uses, controls, and protects.
The software security assurance process begins by identifying and categorizing the information that is to be contained in, or used by, the software. The information should be categorized according to its sensitivity. For example, in the lowest category, the impact of a security violation is minimal (i.e. the impact on the software owner's mission, functions, or reputation is negligible). For a top category, however, the impact may pose a threat to human life; may have an irreparable impact on software owner's missions, functions, image, or reputation; or may result in the loss of significant assets or resources.
Once the information is categorized, security requirements can be developed. The security requirements should address access control, including network access and physical access; data management and data access; environmental controls (power, air conditioning, etc.) and off-line storage; human resource security; and audit trails and usage records.
What causes software security problems?
All security vulnerabilities in software are the result of security bugs, or defects, within the software. In most cases, these defects are created by two primary causes: (1) non-conformance, or a failure to satisfy requirements; and (2) an error or omission in the software requirements.
Non-conformance, or a failure to satisfy requirements
A non-conformance may be simple–the most common is a coding error or defect–or more complex (i.e., a subtle timing error or input validation error). The important point about non-conformance is that verification and validation techniques are designed to detect them and security assurance techniques are designed to prevent them. Improvements in these methods, through a software security assurance program, can improve the security of software.
Errors or omissions in software requirements
The most serious security problems with software-based systems are those that develop when the software requirements are incorrect, inappropriate, or incomplete for the system situation. Unfortunately, errors or omissions in requirements are more difficult to identify. For example, the software may perform exactly as required under normal use, but the requirements may not correctly deal with some system state. When the system enters this problem state, unexpected and undesirable behavior may result. This type of problem cannot be handled within the software discipline; it results from a failure of the system and software engineering processes which developed and allocated the system requirements to the software.
Software security assurance activities
There are two basic types of Software Security Assurance activities.
Some focus on ensuring that information processed by an information system is assigned a proper sensitivity category, and that the appropriate protection requirements have been developed and met in the system.
Others focus on ensuring the control and protection of the software, as well as that of the software support tools and data.
At a minimum, a software security assurance program should ensure that:
A security evaluation has been performed for the software.
Security requirements have been established for the software.
Security requirements have been established for the software development and/or operations and maintenance (O&M) processes.
Each software review, or audit, includes an evaluation of the security requirements.
A configuration management and corrective action process is in place to provide security for the existing software and to ensure that any proposed changes do not inadvertently create security violations or vulnerabilities.
Physical security for the software is adequate.
Building in security
Improving the software development process and building better software are ways to improve software security, by producing software with fewer defects and vulnerabilities. A first-order approach is to identify the critical software components that control security-related functions and pay special attention to them throughout the development and testing process. This approach helps to focus scarce security resources on the most critical areas.
Tools and techniques
There are many commercial off-the-shelf (COTS) software packages that are available to support software security assurance activities. However, before they are used, these tools must be carefully evaluated and their effectiveness must be assured.
Common weaknesses enumeration
One way to improve software security is to gain a better understanding of the most common weaknesses that can affect software security. With that in mind, there is a current community-based program called the Common Weaknesses Enumeration project, which is sponsored by The Mitre Corporation to identify and describe such weaknesses. The list, which is currently in a very preliminary form, contains descriptions of common software weaknesses, faults, and flaws.
Security architecture/design analysis
Security architecture/design analysis verifies that the software design correctly implements security requirements. Generally speaking, there are four basic techniques that are used for security architecture/design analysis.
Logic analysis
Logic analysis evaluates the equations, algorithms, and control logic of the software design.
Data analysis
Data analysis evaluates the description and intended usage of each data item used in design of the software component. The use of interrupts and their effect on data should receive special attention to ensure interrupt handling routines do not alter critical data used by other routines.
Interface analysis
Interface analysis verifies the proper design of a software component's interfaces with other components of the system, including computer hardware, software, and end-users.
Constraint analysis
Constraint analysis evaluates the design of a software component against restrictions imposed by requirements and real-world limitations. The design must be responsive to all known or anticipated restrictions on the software component. These restrictions may include timing, sizing, and throughput constraints, input and output data limitations, equation and algorithm limitations, and other design limitations.
Secure code reviews, inspections, and walkthroughs
Code analysis verifies that the software source code is written correctly, implements the desired design, and does not violate any security requirements. Generally speaking, the techniques used in the performance of code analysis mirror those used in design analysis.
Secure Code reviews are conducted during and at the end of the development phase to determine whether established security requirements, security design concepts, and security-related specifications have been satisfied. These reviews typically consist of the presentation of material to a review group. Secure code reviews are most effective when conducted by personnel who have not been directly involved in the development of the software being reviewed.
Informal reviews
Informal secure code reviews can be conducted on an as-needed basis. To conduct an informal review, the developer simply selects one or more reviewer(s) and provides and/or presents the material to be reviewed. The material may be as informal as pseudo-code or hand-written documentation.
Formal reviews
Formal secure code reviews are conducted at the end of the development phase for each software component. The client of the software appoints the formal review group, who may make or affect a "go/no-go" decision to proceed to the next step of the software development life cycle.
Inspections and walkthroughs
A secure code inspection or walkthrough is a detailed examination of a product on a step-by-step or line-by-line (of source code) basis. The purpose of conducting secure code inspections or walkthroughs is to find errors. Typically, the group that does an inspection or walkthrough is composed of peers from development, security engineering and quality assurance.
Security testing
Software security testing, which includes penetration testing, confirms the results of design and code analysis, investigates software behaviour, and verifies that the software complies with security requirements. Special security testing, conducted in accordance with a security test plan and procedures, establishes the compliance of the software with the security requirements. Security testing focuses on locating software weaknesses and identifying extreme or unexpected situations that could cause the software to fail in ways that would cause a violation of security requirements. Security testing efforts are often limited to the software requirements that are classified as "critical" security items.
See also
Secure by design
Computer security
Security engineering
Software protection
References
Security engineering
Software quality | Software security assurance | [
"Engineering"
] | 1,688 | [
"Systems engineering",
"Security engineering"
] |
5,655,393 | https://en.wikipedia.org/wiki/Mineral%20hydration | In inorganic chemistry, mineral hydration is a reaction which adds water to the crystal structure of a mineral, usually creating a new mineral, commonly called a hydrate.
In geological terms, the process of mineral hydration is known as retrograde alteration and is a process occurring in retrograde metamorphism. It commonly accompanies metasomatism and is often a feature of wall rock alteration around ore bodies. Hydration of minerals occurs generally in concert with hydrothermal circulation which may be driven by tectonic or igneous activity.
Processes
There are two main ways in which minerals hydrate. One is conversion of an oxide to a double hydroxide, as with the hydration of calcium oxide—CaO—to calcium hydroxide—Ca(OH)2. The other is with the incorporation of water molecules directly into the crystalline structure of a new mineral, as with the hydration of feldspars to clay minerals, garnet to chlorite, or kyanite to muscovite.
Mineral hydration is also a process in the regolith that results in conversion of silicate minerals into clay minerals.
Some mineral structures, for example, montmorillonite, are capable of including a variable amount of water without significant change to the mineral structure.
Hydration is the mechanism by which hydraulic binders such as Portland cement develop strength. A hydraulic binder is a material that can set and harden submerged in water by forming insoluble products in a hydration reaction. The term hydraulicity or hydraulic activity is indicative of the chemical affinity of the hydration reaction.
Examples of hydrated minerals
Examples of hydrated minerals include:
silicates (, )
phyllosilicates, clay minerals "commonly found on Earth as weathering products of rocks or in hydrothermal systems"
chlorite
muscovite
non-silicates
oxides (, , , etc.) and oxy-hydroxides
brucite,
goethite, FeO(OH)
carbonates (, etc.)
hydromagnesite,
ikaite, , the unstable hexahydrate form of calcium carbonate
hydroxylated minerals
saponite
talc
hydroxysulfides (mixed sulfides-hydroxides)
tochilinite, a hydroxysulfide or hydrated sulfide mineral of iron(II) and magnesium of chemical formula: , also written , in IMA notation
valleriite, an uncommon sulfide-hydroxide mineral of iron(II) and copper of chemical formula: , or
See also
Clay-water interaction
Hydration reaction
Iron(III) oxide-hydroxide
Ferrihydrite
References
Metamorphic petrology
Inorganic reactions
es:Hidratación mineral | Mineral hydration | [
"Chemistry"
] | 554 | [
"Hydrate minerals",
"Inorganic reactions",
"Hydrates"
] |
5,655,436 | https://en.wikipedia.org/wiki/Antibody%20microarray | An antibody microarray (also known as antibody array) is a specific form of protein microarray. In this technology, a collection of captured antibodies are spotted and fixed on a solid surface such as glass, plastic, membrane, or silicon chip, and the interaction between the antibody and its target antigen is detected. Antibody microarrays are often used for detecting protein expression from various biofluids including serum, plasma and cell or tissue lysates. Antibody arrays may be used for both basic research and medical and diagnostic applications.
Background
The concept and methodology of antibody microarrays were first introduced by Tse Wen Chang in 1983 in a scientific publication and a series of patents, when he was working at Centocor in Malvern, Pennsylvania. Chang coined the term “antibody matrix” and discussed “array” arrangement of minute antibody spots on small glass or plastic surfaces. He demonstrated that a 10×10 (100 in total) and 20×20 (400 in total) grid of antibody spots could be placed on a 1×1 cm surface. He also estimated that if an antibody is coated at a 10 μg/mL concentration, which is optimal for most antibodies, 1 mg of antibody can make 2,000,000 dots of 0.25 mm diameter. Chang's invention focused on the employment of antibody microarrays for the detection and quantification of cells bearing certain surface antigens, such as CD antigens and HLA allotypic antigens, particulate antigens, such as viruses and bacteria, and soluble antigens. The principle of "one sample application, multiple determinations", assay configuration, and mechanics for placing absorbent dots described in the paper and patents should be generally applicable to different kinds of microarrays. When Tse Wen Chang and Nancy T. Chang were setting up Tanox, Inc. in Houston, Texas in 1986, they purchased the rights on the antibody matrix patents from Centocor as part of the technology base to build their new startup. Their first product in development was an assay, termed “immunosorbent cytometry”, which could be employed to monitor the immune status, i.e., the concentrations and ratios of CD3+, CD4+, and CD8+ T cells, in the blood of HIV-infected individuals.
The theoretical background for protein microarray-based ligand binding assays was further developed by Roger Ekins and colleagues in the late 1980s. According to the model, antibody microarrays would not only permit simultaneous screening of an analyte panel, but would also be more sensitive and rapid than conventional screening methods. Interest in screening large protein sets only arose as a result of the achievements in genomics by DNA microarrays and the Human Genome Project.
The first array approaches attempted to miniaturize biochemical and immunobiological assays usually performed in 96-well microtiter plates. While 96-well plate-based antibody arrays have high-throughput capability, the small surface area in each well limits the number of antibody spots and thus, the number of analytes detected. Other solid supports, such as glass slides and nitrocellulose membranes, were subsequently utilized to develop arrays which could accommodate larger panels of antibodies. Nitrocellulose membrane-based arrays are flexible, easy to handle, and have increased protein binding capacity, but are less amenable to high throughput or automated processing. Chemically derivatized glass slides allow for printing of sub-microliter sized antibody spots, reducing the array surface area without sacrificing spot density. This in turn reduces the volume of sample consumed. Glass slide-based arrays, owing to their smooth and rigid structure, can also be easily fitted to high-throughput liquid handling systems.
Most antibody array systems employ 1 of 2 non-competitive methods of immunodetection: single-antibody (label-based) detection and 2-antibody (sandwich-based) detection. The latter method, in which analyte detection requires the binding of 2 distinct antibodies (a capture antibody and a reporter antibody, each binding to a unique epitope), confers greater specificity and lower background signal compared with label-based immunodetection (where only 1 capture antibody is used and detection is achieved by chemically labeling all proteins in the starting sample). Sandwich-based antibody arrays usually attain the highest specificity and sensitivity (ng – pg levels) of any array format; their reproducibility also enables quantitative analysis to be performed. Due to the difficulty of developing matched antibody pairs that are compatible with all other antibodies in the panel, small arrays often make use of a sandwich approach. Conversely, high-density arrays are easier to develop at a lower cost using the single antibody label-based approach. In this methodology, one set of specific antibodies is used and all the proteins in a sample are labelled directly by fluorescent dyes or haptens.
Initial uses of antibody-based array systems included detecting IgGs and specific subclasses, analyzing antigens, screening recombinant antibodies, studying yeast protein kinases, analyzing autoimmune antibodies, and examining protein-protein interactions. The first approach to simultaneously detect multiple cytokines from physiological samples using antibody array technology was by Ruo-Pan Huang and colleagues in 2001. Their approach used Hybond ECL membranes to detect a small panel of 24 cytokines from cell culture conditioned media and patient's sera and was able to profile cytokine expression at physiological levels. Huang took this technology and started a new business, RayBiotech, Inc., the first to successfully commercialize a planar antibody array.
In the last ten years, the sensitivity of the method was improved by an optimization of the surface chemistry as well as dedicated protocols for their chemical labeling. Currently, the sensitivity of antibody arrays is comparable to that of ELISA and antibody arrays are regularly used for profiling experiments on tissue samples, plasma or serum samples and many other sample types. One main focus in antibody array based profiling studies is biomarker discovery, specifically for cancer. For cancer-related research, the development and application of an antibody array comprising 810 different cancer-related antibodies was reported in 2010. Also in 2010, an antibody array comprising 507 cytokines, chemokines, adipokines, growth factors, angiogenic factors, proteases, soluble receptors, soluble adhesion molecules, and other proteins was used to screen the serum of ovarian cancer patients and healthy individuals and found a significant difference in protein expression between normal and cancer samples. More recently, antibody arrays have helped determine specific allergy-related serum proteins whose levels are associated with glioma and can reduce the risk years before diagnosis. Protein profiling with antibody arrays have also proven successful in areas other than cancer research, specifically in neurological diseases such as Alzheimer's. A number of studies have attempted to identify biomarker panels that can distinguish Alzheimer's patients, and many have used antibody arrays in this process. Jaeger and colleagues measured nearly 600 circulatory proteins to discover biological pathways and networks affected in Alzheimer's and explored the positive and negative relationships of the levels of those individual proteins and networks with the cognitive performance of Alzheimer's patients. Currently the largest commercially available sandwich-based antibody array detects 1000 different proteins. In addition, antibody microarray based protein profiling services are available analyzing protein abundance and protein phosphorylation or ubiquitinylation status of 1030 proteins in parallel.
Antibody arrays are often used for detecting protein expression from many sample types, but also in those with various preparations. Jiang and colleagues illustrated nicely the correlation between array protein expression in two different blood preparations: serum and dried blood spots. These different blood sample preparations were analyzed using three antibody array platforms: sandwich-based, quantitative, and label-based, and a strong correlation in protein expression was found, suggesting that dried blood spots, which are a more convenient, safe, and inexpensive means of obtaining blood especially in non-hospitalized public health areas, can be used effectively with antibody array analysis for biomarker discovery, protein profiling, and disease screening, diagnosis, and treatment.
Applications
Using antibody microarray in different medical diagnostic areas has attracted researchers attention. Digital bioassay is an example of such research domains. In this technology, an array of microwells on a glass/polymer chip are seeded with magnetic beads (coated with fluorescent tagged antibodies), subjected to targeted antigens and then characterised by a microscope through counting fluorescing wells. A cost-effective fabrication platform (using OSTE polymers) for such microwell arrays has been recently demonstrated and the bio-assay model system has been successfully characterised. Furthermore, immunoassays on thiol-ene "synthetic paper" micropillar scaffolds have shown to generate a superior fluorescence signal.
See also
ELISA
Protein microarray
DNA microarray
Tissue microarray
Chemical compound microarray
Microarray imprinting and surface energy patterning
References
Microarrays
1983 introductions
Reagents for biochemistry | Antibody microarray | [
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5,655,722 | https://en.wikipedia.org/wiki/Aquadag | Aquadag is a trade name for a water-based colloidal graphite coating commonly used in cathode ray tubes (CRTs). It is manufactured by Acheson Industries, a subsidiary of ICI. The name is a shortened form of "Aqueous Deflocculated Acheson Graphite", but has become a generic term for conductive graphite coatings used in vacuum tubes. Other related products include Oildag, Electrodag and Molydag. Deflocculation refers to the distribution of powdered high purity graphite in an aqueous solution containing approximately 2% to 10% by weight of various Tannic/Gallotannic acid variants and separating the colloidal graphite suspension from the remaining unsuspended graphite particulates. The product names are often printed with DAG in upper case (e.g. AquaDAG). It is used as an electrically conductive coating on insulating surfaces, and as a lubricant.
Properties
Aquadag consists of a dispersion of colloidal graphite in distilled water. It is provided in concentrated paste form and is usually diluted with distilled water to a desired consistency before application. It can be applied by brushing, swabbing, spraying, or dipping, after which the surface is dried, leaving a layer of pure graphite.
After drying the coating is electrically conductive. Its resistance and other electrical properties vary with degree of dilution and application method. When diluted 1:1 and applied by brush its resistance is:
Air-dried ~800 ohms per square
Heated to 200 °C ~500 ohms per square
Heated to 300 °C ~20–30 ohms per square
Use in cathode ray tubes
A conductive aquadag coating applied to the inside of the glass envelope of cathode ray tubes, serves as a high-voltage electrode. The coating covers the inside walls of the "bell" of the CRT tube, from just inside the neck, and stops just short of the screen. Due to the graphite, it is electrically conductive and forms part of the high-voltage positive electrode, the second anode, which accelerates the electron beam. The second anode is a metal cylinder inside the neck of the tube, connected to a high positive voltage of 18 to 25 kilovolts. It has spring clips, which press against the walls of the tube, making contact with the aquadag coating so it also carries this high positive voltage. The electron beam from the electron gun in the neck of the tube is accelerated by the high voltage on the anode and passes through it to strike the screen.
The aquadag coating has two functions: it maintains a uniform electric field inside the tube near the screen, so the electron beam remains collimated and is not distorted by external fields, and it collects the electrons after they have hit the screen, serving as the return path for the cathode current. When the electron beam hits the screen, in addition to causing the fluorescent phosphor coating to give off light, it also knocks other electrons out of the surface. These secondary electrons are attracted to the high positive voltage of the coating and return through it to the anode power supply. Without the coating a negative space charge would develop near the screen, deflecting the electron beam. A typical value of beam current collected by the anode coating is 0.6 mA.
In some CRTs the aquadag coating performs a third function, as a filter capacitor for the high-voltage anode supply. A second conductive coating is applied to part of the outside of the tube facing the inside coating. This outside coating is connected to the ground side of the anode supply, thus the full anode voltage is applied between the coatings. The sandwich of the two coatings separated by the dielectric glass wall of the tube form a final capacitor to filter out ripple from the anode supply. Although the capacitance is small, around 500 pF, due to the low anode current it is sufficient to act as a filter capacitor.
In the television tube manufacturing industry, the manufacturing step that applies the aquadag is called "dagging".
Other uses
Aside from its use in the production of CRTs, Aquadag is used in many types of high-voltage lab apparatus where a conductive coating is needed on an insulating surface. The surfaces of some metals (most notably aluminum) can develop nonconductive oxide layers, which tend to disrupt the electrostatic field produced around the surface of the metal when used as an electrode. Aquadag is not subject to such effects and provides a completely uniform equipotential surface for electrostatics.
Producers of continuous filament fiberglass will coat their product with Aquadag when a conductive property is required.
Aquadag was also used in the production of some copper oxide rectifiers, to help make the ohmic connections to their counterelectrodes.
Other dags
There are also deflocculated graphite products dispersed in liquids other than water. Acheson has extended the use of the dag brandname to non-graphite products e.g. the copper-based Electrodag 437 conductive paint.
References
External links
AquaDAG Product data sheet from
Vacuum tubes | Aquadag | [
"Physics"
] | 1,103 | [
"Vacuum tubes",
"Vacuum",
"Matter"
] |
5,656,456 | https://en.wikipedia.org/wiki/Java%20Analysis%20Studio | Java Analysis Studio (JAS) is an object oriented data analysis package developed for the analysis of particle physics data. The latest major version is JAS3.
JAS3 is a fully AIDA-compliant data analysis system. It is popular for data analysis in areas of particle physics which are familiar with the Java programming language.
The Studio uses many other libraries from the FreeHEP project.
External links
Java Analysis Studio 3 website
AIDA: Abstract Interfaces for Data Analysis — open interfaces and formats for particle physics data processing
Data analysis software
Experimental particle physics
Free software programmed in Java (programming language)
Free statistical software
Numerical software
Physics software | Java Analysis Studio | [
"Physics",
"Mathematics"
] | 128 | [
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"Computational physics",
"Experimental physics",
"Particle physics",
"Numerical software",
"Experimental particle physics",
"Particle physics stubs",
"Computational physics stubs",
"Physics software"
] |
11,048,177 | https://en.wikipedia.org/wiki/Ellerman%20bomb | In solar physics, Ellerman bombs are intense, small-scale brightenings in the Sun's photosphere. They are only observed in the wings of the Hα, Hβ, and Hγ hydrogen spectral lines and take place in emerging flux regions where emerging magnetic fields interact with existing fields. They are named after Ferdinand Ellerman who studied them in detail in the 20th century.
History
Intense brightenings resembling what would later be referred to as Ellerman bombs were first reported by Walter M. Mitchell in 1909. In 1917, observations of this phenomenon made at the Mount Wilson Solar Observatory were described in detail by Ferdinand Ellerman. He referred to them as "solar hydrogen bombs" in reference to the phenomenon only appearing in observations of hydrogen spectral lines.
Description
As originally described in Ellerman's 1917 paper, Ellerman bombs are intense brightenings in the wings of the Hα, Hβ, and Hγ hydrogen spectral lines with no brightening of the line cores or of other spectral lines. They occur in intergranular lanes in the photosphere exclusively at the sites of emerging flux regions where emerging vertical magnetic fields interact with the existing intergranular field. This interaction is suggested to result in magnetic reconnection, producing the brightenings associated with Ellerman bombs.
The lack of observed brightening of the Hα core is attributed to Ellerman bombs being a photospheric phenomenon. In growing active regions, dense chromospheric, Hα fibrils form a canopy above the photosphere blocking Hα emission from Ellerman bombs below. As a result, only emission in the Hα wings pass through and are observed.
References
Solar phenomena
Sun | Ellerman bomb | [
"Physics"
] | 345 | [
"Physical phenomena",
"Stellar phenomena",
"Solar phenomena"
] |
11,050,654 | https://en.wikipedia.org/wiki/Pharmaconomist | In Denmark (including Greenland and Faroe Islands), pharmaconomists () are experts in pharmaceuticals () who have trained with a 3-year tertiary degree. Pharmaconomy () describes either their professional practice or their training courses.
Work
The majority of the Danish pharmaconomists work at community pharmacies (chemists' shops or drug stores) and at hospital pharmacies and hospitals.
Some pharmaconomists work within the chemical industry, the pharmaceutical industry and in medical or clinical laboratories.
Other pharmaconomists teach pharmacy students and pharmaconomy students at colleges or universities, such as at the University of Copenhagen's Faculty of Health and Medical Sciences or at the Pharmakon—Danish College of Pharmacy Practice.
Pharmaconomists are also employed by the Danish Ministry of Interior and Health, Danish Medicines Agency and Danish Association of Pharmacies.
Some pharmaconomists do work as pharmaceutical consultants.
Education
The 3 year higher education corresponds to 180 ECTS points (European Credit Transfer and Accumulation System).
Pharmakon—Danish College of Pharmacy Practice
During his or her education programme at Pharmakon—Danish College of Pharmacy Practice, the pharmaconomist student studies human and animal anatomy, physiology, pathology, pharmacology, pharmaconomy, pharmacy practice, pharmaceutics, toxicology, pharmacognosy, clinical pharmacy, pharmacotherapy, pharmaceutical sciences, chemistry, pharmaceutical chemistry, biochemistry, biology, microbiology, molecular biology, genetics, cytology, medicine, veterinary medicine, zoology, diagnosis, medical prescription, pharmacy law, medical sociology, patient safety, health care, psychology, psychiatry, pedagogy, communication, information technology (IT), bioethics, medical ethics, safety, leadership, organization, logistics, economy, quality assurance (QA), sales and marketing.
Difference between a pharmaconomist and a pharmacist
There are two different professional groups with pharmaceutical education in Denmark:
Pharmaconomists (with a 3-year higher tertiary education)
Pharmacists (with a 5-year higher tertiary education)
Due to his or her higher education as a health professional, the pharmaconomist has by law the same independent competence in all Danish pharmacies as a pharmacist — i.e. for example to dispense and check medical prescriptions, to counsel and advise patients/customers about the use of medicine/pharmaceuticals and to dispense, sell and provide information about medical prescriptions and about prescription medicine and over-the-counter medicine (OTC).
The pharmaconomist also undertakes specialist and managerial operation of pharmacies and undertakes managerial duty service.
The only difference by law is that only a pharmacist may own a Danish pharmacy — i.e. become a pharmacy owner.
Like pharmacists, pharmaconomists can work as pharmacy managers and HR managers (or as chief pharmaconomists).
Trade union
The Danish Association of Pharmaconomists is a trade union who represents about 5,700 pharmaconomists in Denmark (i.e. 98% of all Danish pharmaconomists).
Translation into other languages
The Danish title farmakonom (pharmaconomist) comes from the Greek "pharmakon" (meaning "pharmaceuticals") and "nom" (meaning "expert in").
In Denmark a pharmaconomist is also referred to as lægemiddelkyndig (expert in pharmaceuticals). Lægemiddelkyndig comes from the Danish "lægemiddel" (meaning "pharmaceuticals") and "kyndig" (meaning "expert in").
The title "pharmaconomist" in other languages:
English: pharmaconomist (plural: pharmaconomists)
Danish: farmakonom (plural: farmakonomer)
Faroese: farmakonomur (plural: farmakonomar)
French: pharmaconome (plural: pharmaconomes)
German: Pharmakonom (plural: Pharmakonomen)
Greenlandic: farmakonomit (plural: farmakonominullu)
Italian: farmaconomista (plural: farmaconomisti)
Spanish: farmaconomista (plural: farmaconomistas)
Swedish: farmakonom (plural: farmakonomer)
The title "expert in pharmaceuticals" in other languages:
English: expert in pharmaceuticals (plural: experts in pharmaceuticals)
Danish: lægemiddelkyndig (plural: lægemiddelkyndige)
French: expert en medicaments (plural: experts en medicaments)
German: Arzneimittelexperte (plural: Arzneimittelexperten)
Italian: esperto in farmaci (plural: esperti in farmaci)
Spanish: experto en fármacos (plural: expertos en fármacos)
Swedish: läkemedelsexpert (plural: läkemedelsexperter)
The academic discipline of "pharmaconomy" in other languages:
English: pharmaconomy
Danish: farmakonomi
German: Pharmakonomie
French: pharmaconomie
Italian: farmaconomia
Spanish: farmaconomía
Swedish: farmakonomi
See also
Professional Further Education in Clinical Pharmacy and Public Health
History of pharmacy
Sources
Pharmakon—Danish College of Pharmacy Practice
The Danish Association of Pharmaconomists
The Danish Pharmaceutical Association
Official Curriculum of the Danish Education of Pharmaconomists (September 2007)
Official Executive Order on the Education of Pharmaconomists (June 2007)
Information about pharmaconomists
Pharmaceutical industry
Education in Denmark
Pharmacy in Denmark
Pharmaconomists | Pharmaconomist | [
"Chemistry",
"Biology"
] | 1,256 | [
"Pharmaceutical industry",
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"Life sciences industry"
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11,051,720 | https://en.wikipedia.org/wiki/Solar-powered%20desalination%20unit | A solar-powered desalination unit produces potable water from saline water through direct or indirect methods of desalination powered by sunlight. Solar energy is the most promising renewable energy source due to its ability to drive the more popular thermal desalination systems directly through solar collectors and to drive physical and chemical desalination systems indirectly through photovoltaic cells.
Direct solar desalination produces distillate directly in the solar collector. An example would be a solar still which traps the Sun's energy to obtain freshwater through the process of evaporation and condensation. Indirect solar desalination incorporates solar energy collection systems with conventional desalination systems such as multi-stage flash distillation, multiple effect evaporation, freeze separation or reverse osmosis to produce freshwater.
Direct solar desalination
Solar stills
One type of solar desalination unit is a solar still, it is also similar to a condensation trap. A solar still is a simple way of distilling water, using the heat of the Sun to drive evaporation from humid soil, and ambient air to cool a condenser film. Two basic types of solar stills are box and pit stills. In a pit still, impure water is contained outside the collector, where it is evaporated by sunlight shining through clear plastic. The pure water vapor condenses on the cool inside plastic surface and drips down from the weighted low point, where it is collected and removed. The box type is more sophisticated. The basic principles of solar water distillation are simple, yet effective, as distillation replicates the way nature makes rain. The sun's energy heats water to the point of evaporation. As the water evaporates, water vapor rises, condensing on the glass surface for collection. This process removes impurities, such as salts and heavy metals, and eliminates microbiological organisms. The end result is water cleaner than the purest rainwater.
Indirect solar desalination
Indirect solar desalination systems comprise two sub-systems: a solar collection system and a desalination system. The solar collection system is used, either to collect heat using solar collectors and supply it via a heat exchanger to a thermal desalination process, or to convert electromagnetic solar radiation to electricity using photovoltaic cells to power an electricity-driven desalination process.
Solar-powered reverse osmosis
Osmosis is a natural phenomenon in which water passes through a membrane from a lower to a higher concentration solution. The flow of water can be reversed if a pressure larger than the osmotic pressure is applied on the higher concentration side. In Reverse osmosis desalination systems, seawater pressure is raised above the natural osmotic pressure, forcing pure water through membrane pores to the fresh water side. Reverse osmosis (RO) is the most common desalination process in terms of installed capacity due to its superior energy efficiency compared to thermal desalination systems, despite requiring extensive water pre-treatment. Furthermore, part of the consumed mechanical energy can be reclaimed from the concentrated brine effluent with an energy recovery device.
Solar-powered RO desalination is common in demonstration plants due to the modularity and scalability of both photovoltaic (PV) and RO systems. A detailed economic analysis and a thorough optimisation strategy of PV powered RO desalination were carried out with favorable results reported. Economic and reliability considerations are the main challenges to improving PV powered RO desalination systems. However, the quickly dropping PV panel costs are making solar-powered desalination ever more feasible.
A solar powered desalination unit designed for remote communities has been tested in the Northern Territory of Australia. The "reverse-osmosis solar installation" (ROSI) uses membrane filtration to provide a reliable and clean drinking water stream from sources such as brackish groundwater. Solar energy overcomes the usually high-energy operating costs as well as greenhouse emissions of conventional reverse osmosis systems. ROSI can also remove trace contaminants such as arsenic and uranium that may cause certain health problems, and minerals such as calcium carbonate which causes water hardness.
Project leader Dr Andrea Schaefer from the University of Wollongong's Faculty of Engineering said ROSI has the potential to bring clean water to remote communities throughout Australia that do not have access to a town water supply and/or the electricity grid.
Groundwater (which may contain dissolved salts or other contaminants) or surface water (which may have high turbidity or contain microorganisms) is pumped into a tank with an ultrafiltration membrane, which removes viruses and bacteria. This water is fit for cleaning and bathing. Ten percent of that water undergoes nanofiltration and reverse osmosis in the second stage of purification, which removes salts and trace contaminants, producing drinking water. A photovoltaic solar array tracks the Sun and powers the pumps needed to process the water, using the plentiful sunlight available in remote regions of Australia not served by the power grid.
Solar photo voltaic power is considered a viable option to power a reverse osmosis desalination plant. The techno-economics both in standalone mode and in PV-biodisel hybrid mode for capacities from 0.05 MLD to 300 MLD were examined by researchers at IIT Madras. As a technology demonstrator, a plant of 500 litre /day capacity has been designed, installed and functional there.
Energy storage
While the intermittent nature of sunlight and its variable intensity throughout the day makes desalination during nighttime challenging, several energy storage options can be used to permit 24 hour operation. Batteries can store solar energy for use at night. Thermal energy storage systems ensure constant performance at night or on cloudy days, improving overall efficiency. Alternatively, stored gravitational energy can be harnessed to provide energy to a solar-powered reverse osmosis unit during non-sunlight hours.
See also
Solar desalination
Solar still
References
Solar energy in Australia
Appropriate technology
Water treatment
Solar-powered devices
Membrane technology | Solar-powered desalination unit | [
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11,054,533 | https://en.wikipedia.org/wiki/Archard%20equation | The Archard wear equation is a simple model used to describe sliding wear and is based on the theory of asperity contact. The Archard equation was developed much later than (sometimes also known as energy dissipative hypothesis), though both came to the same physical conclusions, that the volume of the removed debris due to wear is proportional to the work done by friction forces. Theodor Reye's model became popular in Europe and it is still taught in university courses of applied mechanics. Until recently, Reye's theory of 1860 has, however, been totally ignored in English and American literature where subsequent works by Ragnar Holm and John Frederick Archard are usually cited. In 1960, and Mikhail Alekseevich Babichev published a similar model as well. In modern literature, the relation is therefore also known as Reye–Archard–Khrushchov wear law. In 2022, the steady-state Archard wear equation was extended into the running-in regime using the bearing ratio curve representing the initial surface topography.
Equation
where:
Q is the total volume of wear debris produced
K is a dimensionless constant
W is the total normal load
L is the sliding distance
H is the hardness of the softest contacting surfaces
Note that is proportional to the work done by the friction forces as described by Reye's hypothesis.
Also, K is obtained from experimental results and depends on several parameters. Among them are surface quality, chemical affinity between the material of two surfaces, surface hardness process, heat transfer between two surfaces and others.
Derivation
The equation can be derived by first examining the behavior of a single asperity.
The local load , supported by an asperity, assumed to have a circular cross-section with a radius , is:
where P is the yield pressure for the asperity, assumed to be deforming plastically. P will be close to the indentation hardness, H, of the asperity.
If the volume of wear debris, , for a particular asperity is a hemisphere sheared off from the asperity, it follows that:
This fragment is formed by the material having slid a distance 2a
Hence, , the wear volume of material produced from this asperity per unit distance moved is:
making the approximation that
However, not all asperities will have had material removed when sliding distance 2a. Therefore, the total wear debris produced per unit distance moved, will be lower than the ratio of W to 3H. This is accounted for by the addition of a dimensionless constant K, which also incorporates the factor 3 above. These operations produce the Archard equation as given above. Archard interpreted K factor as a probability of forming wear debris from asperity encounters. Typically for 'mild' wear, K ≈ 10−8, whereas for 'severe' wear, K ≈ 10−2. Recently, it has been shown that there exists a critical length scale that controls the wear debris formation at the asperity level. This length scale defines a critical junction size, where bigger junctions produce debris, while smaller ones deform plastically.
See also
References
Further reading
https://patents.google.com/patent/DE102005060024A1/de (Mentions the term "Reye-Hypothese")
Surfaces
Materials science
Equations
Tribology | Archard equation | [
"Physics",
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11,059,015 | https://en.wikipedia.org/wiki/Steam%20crane | A steam crane is a crane powered by a steam engine. It may be fixed or mobile and, if mobile, it may run on rail tracks, caterpillar tracks, road wheels, or be mounted on a barge. It usually has a vertical boiler placed at the back so that the weight of the boiler counterbalances the weight of the jib and load.
They were very common as railway breakdown cranes, and several have been preserved on heritage railways in the United Kingdom.
Manufacturers
Black Hawthorn of Gateshead (unrestored example at Beamish Museum
Joseph Booth & Bros of Leeds
Coles Cranes of Derby (restored example at Beamish Museum)
Cowans, Sheldon & Company of Carlisle (rail cranes)
Craven Brothers
William Fairbairn & Sons of Manchester
Ransomes & Rapier of Ipswich
Ruston Proctor of Lincoln
Stothert & Pitt of Bath
Thomas Smith & Sons (Rodley) Ltd. of Leeds
See also
Crane
Crane (railroad)
Crane tank
Fairbairn steam crane
Steam engine
Steam shovel
References
External links
Steam cranes inc. Ransomes & Rapier
Cowans & Sheldon steam crane
Nine Elms steam crane
Ransomes & Rapier wartime-ordered 45-ton Steam Breakdown Cranes
Cowans Sheldon 15-ton Steam Cranes
Model steam crane
Cranes (machines)
Crane | Steam crane | [
"Engineering"
] | 255 | [
"Engineering vehicles",
"Mechanical engineering stubs",
"Cranes (machines)",
"Mechanical engineering"
] |
2,254,244 | https://en.wikipedia.org/wiki/Dogbane | Dogbane, dog-bane, dog's bane, and other variations, some of them regional and some transient, are names for certain plants that are reputed to kill or repel dogs; "bane" originally meant "slayer", and was later applied to plants to indicate that they were poisonous to particular creatures.
History of the term
The earliest reference to such names in common English usage was in the 16th century, in which they were applied to various plants in the Apocynaceae, in particular Apocynum. Some plants in the Asclepiadoideae, now a subfamily of the Apocynaceae, but until recently regarded as the separate family Asclepiadaceae, were also called dogbane even before the two families were united. It is not clear how much earlier the name had been in use in the English language, which originated about 1000 years earlier in mediaeval times. However, centuries before the appearance of the English language, Pedanius Dioscorides, in his De Materia Medica, had already described members of the Apocynaceae, such as Apocynum and Cynanchum by names equivalent to "dogbane"; Apocynum literally means "dog killer" or "dog remover", and "Cynanchum" means "dog strangler". In modern times some species of Nerium, Periploca and Trachelospermum, also in the Apocynaceae, are called dogbane or variants such as "climbing dogbane".
Modern significance of the term "dogbane family"
Some modern sources note "dogbane" as strictly being the species known as 'Indian hemp', Apocynum, though it is doubtful that such a narrow definition could be justified, even if it were enforceable. Still others consider Asclepias (milkweeds) to be the "true" dogbanes; however, when the majority of authors, horticulturists or gardeners refer to the "dogbanes", they are generally always referencing the entire Apocynaceae family, as a whole.
"Dogbane" as a term outside the family Apocynaceae
Common names, either informal or vernacular, are seldom definitive, let alone stable. Some poisonous or offensive plants in practically unrelated families had similar common names in the vernacular and writings of various times; for example an edition De Materia Medica, apparently of the early sixteenth century, mentions that species of Aconitum (family Ranunculaceae) were known as either "dog killer" (or murderer) or "wolf killer" ("...Sunt qui Cynoctonon: qui Lycoctonon... uocent"). Again, in modern times Isocoma menziesii in the family Asteraceae is known in some regions as dogbane.
Recent aberrant application of the term
The term "dogbane" (as well as "cat-scat")—either out of genuine confusion or as a deliberate sales ploy for gardeners desiring a natural animal repellent—has been used without obvious justification to several other groups of plants, such as some species of Plectranthus (ironically, a genus in the catnip subfamily Nepetoideae of the mint family Lamiaceae). While none have been reported to be especially harmful, or even useful against nuisance animals, in the garden, many—such as Plectranthus (Coleus) caninus—have very fragrant, oily leaves which give an intensely pungent, skunk- or Cannabis-like aroma when brushed, disturbed or touched. At times, simply the wind blowing can trigger the release of the essential oils into the surrounding area. The smell has been reported, by some sources, to keep nuisance animals at bay; however, if a plant is not poisonous or otherwise offensive to them, many animals quickly become accustomed to various botanical aromas and remain unbothered. Oftentimes, these plants are more effective at repelling humans from a given area, as the essential oils are strong, sticky, and exude a distinct aroma of marijuana or skunk-spray, which may linger for hours on the skin, gloves, clothing, or any other surface it contacts.
References
Apocynaceae
Plant common names | Dogbane | [
"Biology"
] | 894 | [
"Plant common names",
"Common names of organisms",
"Plants"
] |
2,255,444 | https://en.wikipedia.org/wiki/Retention%20basin | A retention basin, sometimes called a retention pond, wet detention basin, or storm water management pond (SWMP), is an artificial pond with vegetation around the perimeter and a permanent pool of water in its design. It is used to manage stormwater runoff, for protection against flooding, for erosion control, and to serve as an artificial wetland and improve the water quality in adjacent bodies of water.
It is distinguished from a detention basin, sometimes called a "dry pond", which temporarily stores water after a storm, but eventually empties out at a controlled rate to a downstream water body. It also differs from an infiltration basin which is designed to direct stormwater to groundwater through permeable soils.
Wet ponds are frequently used for water quality improvement, groundwater recharge, flood protection, aesthetic improvement, or any combination of these. Sometimes they act as a replacement for the natural absorption of a forest or other natural process that was lost when an area is developed. As such, these structures are designed to blend into neighborhoods and viewed as an amenity.
In urban areas, impervious surfaces (roofs, roads) reduce the time spent by rainfall before entering into the stormwater drainage system. If left unchecked, this will cause widespread flooding downstream. The function of a stormwater pond is to contain this surge and release it slowly. This slow release mitigates the size and intensity of storm-induced flooding on downstream receiving waters. Stormwater ponds also collect suspended sediments, which are often found in high concentrations in stormwater water due to upstream construction and sand applications to roadways.
Design features
Storm water is typically channeled to a retention basin through a system of street and/or parking lot storm drains, and a network of drain channels or underground pipes. The basins are designed to allow relatively large flows of water to enter, but discharges to receiving waters are limited by outlet structures that function only during very large storm events.
Retention ponds are often landscaped with a variety of grasses, shrubs, and/or aquatic plants to provide bank stability and aesthetic benefits. Vegetation also provides water quality benefits by removing soluble nutrients through uptake. In some areas the ponds can attract nuisance types of wildlife like ducks or Canada geese, particularly where there is minimal landscaping and grasses are mowed. This reduces the ability of foxes, coyotes, and other predators to approach their prey unseen. Such predators tend to hide in the cattails and other tall, thick grass surrounding natural water features.
Proper depth of retention ponds is important for removal of pollutants and maintenance of fish populations. Urban fishing continues to be one of the fastest growing fishing segments as new suburban neighborhoods are built around these aquatic areas.
Other meanings
A retention basin can also be a part of a nuclear reactor used to contain a core meltdown.
See also
Balancing lake (UK)
Nationwide Urban Runoff Program (NURP) – US stormwater research project
Settling basin – for treating agricultural and industrial wastewater
Stream restoration
Surface runoff
Sustainable drainage system
Urban runoff
Water pollution
Jonenbach flood retention basin
References
External links
Virginia retention basin standards
Detention vs. retention – Harris County, Texas Flood Control District
Stormwater Ecological Enhancement Project – University of Florida
The use of retention ponds in residential settings
International Stormwater BMP Database – Performance Data on Urban Stormwater BMPs
Environmental engineering
Hydraulic engineering
Hydrology
Infrastructure
Stormwater management | Retention basin | [
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2,258,380 | https://en.wikipedia.org/wiki/Nambu%E2%80%93Jona-Lasinio%20model | In quantum field theory, the Nambu–Jona-Lasinio model (or more precisely: the Nambu and Jona-Lasinio model) is a complicated effective theory of nucleons and mesons constructed from interacting Dirac fermions with chiral symmetry, paralleling the construction of Cooper pairs from electrons in the BCS theory of superconductivity. The "complicatedness" of the theory has become more natural as it is now seen as a low-energy approximation of the still more basic theory of quantum chromodynamics, which does not work perturbatively at low energies.
Overview
The model is much inspired by the different field of solid state theory, particularly from the BCS breakthrough of 1957. The model was introduced in a joint article of Yoichiro Nambu (who also contributed essentially to the theory of superconductivity, i.e., by the "Nambu formalism") and Giovanni Jona-Lasinio, published in 1961. A subsequent paper included chiral symmetry breaking, isospin and strangeness.
Around that time, the same model was independently considered by Soviet physicists Valentin Vaks and Anatoly Larkin.
The model is quite technical, although based essentially on symmetry principles. It is an example of the importance of four-fermion interactions and is defined in a spacetime with an even number of dimensions. It is still important and is used primarily as an effective although not rigorous low energy substitute for quantum chromodynamics.
The dynamical creation of a condensate from fermion interactions inspired many theories of the breaking of electroweak symmetry, such as technicolor and the top-quark condensate.
Starting with the one-flavor case first, the Lagrangian density is
or, equivalently,
The terms proportional to are an attractive four-fermion interaction, which parallels the BCS theory phonon exchange interaction.
The global symmetry of the model is U(1)Q×U(1)χ where Q is the ordinary charge of the Dirac fermion and χ is the chiral charge.
is actually an inverse squared mass, which represents short-distance physics or
the strong interaction scale, producing an attractive four-fermion interaction.
There is no bare fermion mass term because of the chiral symmetry. However, there will be a chiral condensate (but no confinement) leading to an effective mass term and a spontaneous symmetry breaking of the chiral symmetry, but not the charge symmetry.
With N flavors and the flavor indices represented by the Latin letters a, b, c, the Lagrangian density becomes
Chiral symmetry forbids a bare mass term, but there may be chiral condensates. The global symmetry here is SU(N)L×SU(N)R× U(1)Q × U(1)χ where SU(N)L×SU(N)R acting upon the left-handed flavors and right-handed flavors respectively is the chiral symmetry (in other words, there is no natural correspondence between the left-handed and the right-handed flavors), U(1)Q is the Dirac charge, which is sometimes called the baryon number and U(1)χ is the axial charge. If a chiral condensate forms, then the chiral symmetry is spontaneously broken into a diagonal subgroup SU(N) since the condensate leads to a pairing of the left-handed and the right-handed flavors. The axial charge is also spontaneously broken.
The broken symmetries lead to massless pseudoscalar bosons which are sometimes called pions. See Goldstone boson.
As mentioned, this model is sometimes used as a phenomenological model of quantum chromodynamics in the chiral limit. However, while it is able to model chiral symmetry breaking and chiral condensates, it does not model confinement. Also, the axial symmetry is broken spontaneously in this model, leading to a massless Goldstone boson unlike QCD, where it is broken anomalously.
Since the Nambu–Jona-Lasinio model is nonrenormalizable in four spacetime dimensions, this theory can only be an effective field theory which needs to be UV completed.
See also
Gross–Neveu model
References
External links
Giovanni Jona-Lasinio and Yoichiro Nambu, Nambu-Jona-Lasinio model, Scholarpedia, 5(12):7487, (2010). doi:10.4249/scholarpedia.7487
Quantum chromodynamics
Superconductivity | Nambu–Jona-Lasinio model | [
"Physics",
"Materials_science",
"Engineering"
] | 970 | [
"Physical quantities",
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"Materials science",
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2,258,865 | https://en.wikipedia.org/wiki/Isotopologue | In chemistry, isotopologues (also spelled isotopologs) are molecules that differ only in their isotopic composition. They have the same chemical formula and bonding arrangement of atoms, but at least one atom has a different number of neutrons than the parent.
An example is water, whose hydrogen-related isotopologues are: "light water" (HOH or ), "semi-heavy water" with the deuterium isotope in equal proportion to protium (HDO or ), "heavy water" with two deuterium atoms ( or ); and "super-heavy water" or tritiated water ( or , as well as and , where some or all of the hydrogen is the radioactive tritium isotope). Oxygen-related isotopologues of water include the commonly available form of heavy-oxygen water () and the more difficult to separate version with the isotope. Both elements may be replaced by isotopes, for example in the doubly labeled water isotopologue . Altogether, there are 9 different stable water isotopologues, and 9 radioactive isotopologues involving tritium, for a total of 18. However only certain ratios are possible in mixture, due to prevalent hydrogen swapping.
The atom(s) of the different isotope may be anywhere in a molecule, so the difference is in the net chemical formula. If a compound has several atoms of the same element, any one of them could be the altered one, and it would still be the same isotopologue. When considering the different locations of the same isotope, the term isotopomer, first proposed by Seeman and Paine in 1992, is used.
Isotopomerism is analogous to constitutional isomerism or stereoisomerism of different elements in a structure. Depending on the formula and the symmetry of the structure, there might be several isotopomers of one isotopologue. For example, ethanol has the molecular formula . Mono-deuterated ethanol, or , is an isotopologue of it. The structural formulas and are two isotopomers of that isotopologue.
Singly substituted isotopologues
Analytical chemistry applications
Singly substituted isotopologues may be used for nuclear magnetic resonance experiments, where deuterated solvents such as deuterated chloroform (CDCl or CHCl) do not interfere with the solutes' H signals, and in investigations of the kinetic isotope effect.
Geochemical applications
In the field of stable isotope geochemistry, isotopologues of simple molecules containing rare heavy isotopes of carbon, oxygen, hydrogen, nitrogen, and sulfur are used to trace equilibrium and kinetic processes in natural environments and in Earth's past.
Doubly substituted isotopologues
Measurement of the abundance of clumped isotopes (doubly substituted isotopologues) of gases has been used in the field of stable isotope geochemistry to trace equilibrium and kinetic processes in the environment inaccessible by analysis of singly substituted isotopologues alone.
Currently measured doubly substituted isotopologues include:
Carbon dioxide: COO
Methane: 13CH3D and 12CH2D2
Oxygen: 18O2 and 17O18O
Nitrogen: 15N2
Nitrous oxide: NNO and NNO
Analytical requirements
Because of the relative rarity of the heavy isotopes of C, H, and O, isotope-ratio mass spectrometry (IRMS) of doubly substituted species requires larger volumes of sample gas and longer analysis times than traditional stable isotope measurements, thereby requiring extremely stable instrumentation. Also, the doubly-substituted isotopologues are often subject to isobaric interferences, as in the methane system where CH and CHD ions interfere with measurement of the CHD and CHD species at mass 18. A measurement of such species requires either very high mass resolving power to separate one isobar from another, or modeling of the contributions of the interfering species to the abundance of the species of interest. These analytical challenges are significant: The first publication precisely measuring doubly substituted isotopologues did not appear until 2004, though singly substituted isotopologues had been measured for decades previously.
As an alternative to more conventional gas source IRMS instruments, tunable diode laser absorption spectroscopy has also emerged as a method to measure doubly substituted species free from isobaric interferences, and has been applied to the methane isotopologue CHD.
Equilibrium fractionation
When a light isotope is replaced with a heavy isotope (e.g., C for C), the bond between the two atoms will vibrate more slowly, thereby lowering the zero-point energy of the bond and acting to stabilize the molecule. An isotopologue with a doubly substituted bond is therefore slightly more thermodynamically stable, which will tend to produce a higher abundance of the doubly substituted (or “clumped”) species than predicted by the statistical abundance of each heavy isotope (known as a stochastic distribution of isotopes). This effect increases in magnitude with decreasing temperature, so the abundance of the clumped species is related to the temperature at which the gas was formed or equilibrated. By measuring the abundance of the clumped species in standard gases formed in equilibrium at known temperatures, the thermometer can be calibrated and applied to samples with unknown abundances.
Kinetic fractionation
The abundances of multiply substituted isotopologues can also be affected by kinetic processes. As for singly substituted isotopologues, departures from thermodynamic equilibrium in a doubly-substituted species can implicate the presence of a particular reaction taking place. Photochemistry occurring in the atmosphere has been shown to alter the abundance of O from equilibrium, as has photosynthesis. Measurements of CHD and CHD can identify microbial processing of methane and have been used to demonstrate the significance of quantum tunneling in the formation of methane, as well as mixing and equilibration of multiple methane reservoirs. Variations in the relative abundances of the two NO isotopologues NNO and {{sup>15}}NNO can distinguish whether NO has been produced by bacterial denitrification or by bacterial nitrification.
Multiple substituted isotopologues
Biochemical applications
Multiple substituted isotopologues may be used for nuclear magnetic resonance or mass spectrometry experiments, where isotopologues are used to elucidate metabolic pathways in a qualitative (detect new pathways) or quantitative (detect quantitative share of a pathway) approach. A popular example in biochemistry is the use of uniform labelled glucose (U-13C glucose), which is metabolized by the organism under investigation (e. g. bacterium, plant, or animal) and whose signatures can later be detected in newly formed amino acid or metabolically cycled products.
Mass spectrometry applications
Resulting from either naturally occurring isotopes or artificial isotopic labeling, isotopologues can be used in various mass spectrometry applications.
Applications of natural isotopologues
The relative mass spectral intensity of natural isotopologues, calculable from the fractional abundances of the constituent elements, is exploited by mass spectrometry practitioners in quantitative analysis and unknown compound identification:
To identify the more likely molecular formulas for an unknown compound based on the matching between the observed isotope abundance pattern in an experiment and the expected isotope abundance patterns for given molecular formulas.
To expand the linear dynamic response range of the mass spectrometer by following multiple isotopologues, with an isotopologue of lower abundance still generating linear response even while the isotopologues of higher abundance giving saturated signals.
Applications of isotope labeling
A compound tagged by replacing specific atoms with the corresponding isotopes can facilitate the following mass spectrometry methods:
Metabolic flux analysis (MFA)
Stable isotopically labeled internal standards for quantitative analysis
See also
Mass (mass spectrometry)
Isotope-ratio mass spectrometry
Isotopomer
Clumped isotopes
Isotopocule
References
External links
Fractional abundance of atmospheric isotopologues, SpectralCalc.com
Isotopes | Isotopologue | [
"Physics",
"Chemistry"
] | 1,688 | [
"Isotopes",
"Nuclear physics"
] |
2,259,059 | https://en.wikipedia.org/wiki/Chronometry | Chronometry or horology () is the science studying the measurement of time and timekeeping. Chronometry enables the establishment of standard measurements of time, which have applications in a broad range of social and scientific areas. Horology usually refers specifically to the study of mechanical timekeeping devices, while chronometry is broader in scope, also including biological behaviours with respect to time (biochronometry), as well as the dating of geological material (geochronometry).
Horology is commonly used specifically with reference to the mechanical instruments created to keep time: clocks, watches, clockwork, sundials, hourglasses, clepsydras, timers, time recorders, marine chronometers, and atomic clocks are all examples of instruments used to measure time. People interested in horology are called horologists. That term is used both by people who deal professionally with timekeeping apparatuses, as well as enthusiasts and scholars of horology. Horology and horologists have numerous organizations, both professional associations and more scholarly societies. The largest horological membership organisation globally is the NAWCC, the National Association of Watch and Clock Collectors, which is US based, but also has local chapters elsewhere.
Records of timekeeping are attested during the Paleolithic, in the form of inscriptions made to mark the passing of lunar cycles and measure years. Written calendars were then invented, followed by mechanical devices. The highest levels of precision are presently achieved by atomic clocks, which are used to track the international standard second.
Etymology
Chronometry is derived from two root words, chronos and metron (χρόνος and μέτρον in Ancient Greek respectively), with rough meanings of "time" and "measure". The combination of the two is taken to mean time measuring.
In the Ancient Greek lexicon, meanings and translations differ depending on the source. Chronos, used in relation to time when in definite periods, and linked to dates in time, chronological accuracy, and sometimes in rare cases, refers to a delay. The length of the time it refers ranges from seconds to seasons of the year to lifetimes, it can also concern periods of time wherein some specific event takes place, or persists, or is delayed.
The root word is correlated with the god Chronos in Ancient Greek mythology, who embodied the image of time, originated from out of the primordial chaos. Known as the one who spins the Zodiac Wheel, further evidence of his connection to the progression of time. However, Ancient Greek makes a distinction between two types of time, chronos, the static and continuing progress of present to future, time in a sequential and chronological sense, and Kairos, a concept based in a more abstract sense, representing the opportune moment for action or change to occur.
Kairos (καιρός) carries little emphasis on precise chronology, instead being used as a time specifically fit for something, or also a period of time characterised by some aspect of crisis, also relating to the endtime. It can as well be seen in the light of an advantage, profit, or fruit of a thing, but has also been represented in apocalyptic feeling, and likewise shown as variable between misfortune and success, being likened to a body part vulnerable due to a gap in armor for Homer, benefit or calamity depending on the perspective. It is also referenced in Christian theology, being used as implication of God's action and judgement in circumstances.
Because of the inherent relation between chronos and kairos, their function the Ancient Greek's portrayal and concept of time, understanding one means understanding the other in part. The implication of chronos, an indifferent disposition and eternal essence lies at the core of the science of chronometry, bias is avoided, and definite measurement is favoured.
Subfields
Biochronometry
Biochronometry (also chronobiology or biological chronometry) is the study of biological behaviours and patterns seen in animals with factors based in time. It can be categorised into Circadian rhythms and Circannual cycles. Examples of these behaviours can be: the relation of daily and seasonal tidal cues to the activity of marine plants and animals, the photosynthetic capacity and phototactic responsiveness in algae, or metabolic temperature compensation in bacteria.
Circadian rhythms of various species can be observed through their gross motor function throughout the course of a day. These patterns are more apparent with the day further categorised into activity and rest times. Investigation into a species is conducted through comparisons of free-running and entrained rhythms, where the former is attained from within the species' natural environment and the latter from a subject that has been taught certain behaviours. Circannual rhythms are alike but pertain to patterns within the scale of a year, patterns like migration, moulting, reproduction, and body weight are common examples, research and investigation are achieved with similar methods to circadian patterns.
Circadian and circannual rhythms can be seen in all organisms, in both single and multi-celled organisms. A sub-branch of biochronometry is microbiochronometry (also chronomicrobiology or microbiological chronometry), and is the examination of behavioural sequences and cycles within micro-organisms. Adapting to circadian and circannual rhythms is an essential evolution for living organisms, these studies, as well as educating on the adaptations of organisms also bring to light certain factors affecting many of species' and organisms' responses, and can also be applied to further understand the overall physiology, this can be for humans as well, examples include: factors of human performance, sleep, metabolism, and disease development, which are all connected to biochronometrical cycles.
Mental chronometry
Mental chronometry (also called cognitive chronometry) studies human information processing mechanisms, namely reaction time and perception. As well as a field of chronometry, it also forms a part of cognitive psychology and its contemporary human information processing approach. Research comprises applications of the chronometric paradigms – many of which are related to classical reaction time paradigms from psychophysiology – through measuring reaction times of subjects with varied methods, and contribute to studies in cognition and action. Reaction time models and the process of expressing the temporostructural organisation of human processing mechanisms have an innate computational essence to them. It has been argued that because of this, conceptual frameworks of cognitive psychology cannot be integrated in their typical fashions.
One common method is the use of event-related potentials (ERPs) in stimulus-response experiments. These are fluctuations of generated transient voltages in neural tissues that occur in response to a stimulus event either immediately before or after. This testing emphasises the mental events' time-course and nature and assists in determining the structural functions in human information processing.
Geochronometry
The dating of geological materials makes up the field of geochronometry, and falls within areas of geochronology and stratigraphy, while differing itself from chronostratigraphy. The geochronometric scale is periodic, its units working in powers of 1000, and is based in units of duration, contrasting with the chronostratigraphic scale. The distinctions between the two scales have caused some confusion – even among academic communities.
Geochronometry deals with calculating a precise date of rock sediments and other geological events, giving an idea as to what the history of various areas is, for example, volcanic and magmatic movements and occurrences can be easily recognised, as well as marine deposits, which can be indicators for marine events and even global environmental changes. This dating can be done in a number of ways. All dependable methods – barring the exceptions of thermoluminescence, radioluminescence and ESR (electron spin resonance) dating – are based in radioactive decay, focusing on the degradation of the radioactive parent nuclide and the corresponding daughter product's growth.
By measuring the daughter isotopes in a specific sample its age can be calculated. The preserved conformity of parent and daughter nuclides provides the basis for the radioactive dating of geochronometry, applying the Rutherford Soddy Law of Radioactivity, specifically using the concept of radioactive transformation in the growth of the daughter nuclide.
Thermoluminescence is an extremely useful concept to apply, being used in a diverse amount of areas in science, dating using thermoluminescence is a cheap and convenient method for geochronometry. Thermoluminescence is the production of light from a heated insulator and semi-conductor, it is occasionally confused with incandescent light emissions of a material, a different process despite the many similarities. However, this only occurs if the material has had previous exposure to and absorption of energy from radiation. Importantly, the light emissions of thermoluminescence cannot be repeated. The entire process, from the material's exposure to radiation would have to be repeated to generate another thermoluminescence emission. The age of a material can be determined by measuring the amount of light given off during the heating process, by means of a phototube, as the emission is proportional to the dose of radiation the material absorbed.
Time metrology
Time metrology or time and frequency metrology is the application of metrology for timekeeping, including frequency stability.
Its main tasks are the realization of the second as the SI unit of measurement for time and the establishment of time standards and frequency standards as well as their dissemination.
History
Early humans would have used their basic senses to perceive the time of day, and relied on their biological sense of time to discern the seasons in order to act accordingly. Their physiological and behavioural seasonal cycles mainly being influenced by a melatonin based photoperiod time measurement biological system – which measures the change in daylight within the annual cycle, giving a sense of the time in the year – and their circannual rhythms, providing an anticipation of environmental events months beforehand to increase chances of survival.
There is debate over when the earliest use of lunar calendars was, and over whether some findings constituted as a lunar calendar. Most related findings and materials from the palaeolithic era are fashioned from bones and stone, with various markings from tools. These markings are thought to not have been the result of marks to represent the lunar cycles but non-notational and irregular engravings, a pattern of latter subsidiary marks that disregard the previous design is indicative of the markings being the use of motifs and ritual marking instead.
However, as humans' focus turned to farming the importance and reliance on understanding the rhythms and cycle of the seasons grew, and the unreliability of lunar phases became problematic. An early human accustomed to the phases of the moon would use them as a rule of thumb, and the potential for weather to interfere with reading the cycle further degraded the reliability. The length of a moon is on average less than our current month, not acting as a dependable alternate, so as years progress the room of error between would grow until some other indicator would give indication.
The Ancient Egyptian calendars were among the first calendars made, and the civil calendar even endured for a long period afterwards, surviving past even its culture's collapse and through the early Christian era. It has been assumed to have been invented near 4231 BC by some, but accurate and exact dating is difficult in its era and the invention has been attributed to 3200 BC, when the first historical king of Egypt, Menes, united Upper and Lower Egypt. It was originally based on cycles and phases of the moon, however, Egyptians later realised the calendar was flawed upon noticing the star Sirius rose before sunrise every 365 days, a year as we know it now, and was remade to consist of twelve months of thirty days, with five epagomenal days. The former is referred to as the Ancient Egyptians' lunar calendar, and the latter the civil calendar.
Early calendars often hold an element of their respective culture's traditions and values, for example, the five day intercalary month of the Ancient Egyptian's civil calendar representing the birthdays of the gods Horus, Isis, Set, Osiris and Nephthys. Maya use of a zero date as well as the Tzolkʼin's connection to their thirteen layers of heaven (the product of it and all the human digits, twenty, making the 260-day year of the year) and the length of time between conception and birth in pregnancy.
Museums and libraries
Europe
There are many horology museums and several specialized libraries devoted to the subject. One example is the Royal Greenwich Observatory, which is also the source of the Prime Meridian and the home of the first marine timekeepers accurate enough to determine longitude (made by John Harrison). Other horological museums in the London area include the Clockmakers' Museum, which re-opened at the Science Museum in October 2015, the horological collections at the British Museum, the Science Museum (London), and the Wallace Collection. The Guildhall Library in London contains an extensive public collection on horology. In Upton, also in the United Kingdom, at the headquarters of the British Horological Institute, there is the Museum of Timekeeping. A more specialised museum of horology in the United Kingdom is the Cuckooland Museum in Cheshire, which hosts the world's largest collection of antique cuckoo clocks.
One of the more comprehensive museums dedicated to horology is the Musée international d'horlogerie, in La Chaux-de-Fonds in Switzerland, which contains a public library of horology. The Musée d'Horlogerie du Locle is smaller but located nearby. Other good horological libraries providing public access are at the Musée international d'horlogerie in Switzerland, at La Chaux-de-Fonds, and at Le Locle.
In France, Besançon has the Musée du Temps (Museum of Time) in the historic Palais Grenvelle. In Serpa and Évora, in Portugal, there is the Museu do Relógio. In Germany, there is the Deutsches Uhrenmuseum in Furtwangen im Schwarzwald, in the Black Forest, which contains a public library of horology.
North America
The two leading specialised horological museums in North America are the National Watch and Clock Museum in Columbia, Pennsylvania, and the American Clock and Watch Museum in Bristol, Connecticut. Another museum dedicated to clocks is the Willard House and Clock Museum in Grafton, Massachusetts. One of the most comprehensive horological libraries open to the public is the National Watch and Clock Library in Columbia, Pennsylvania.
Organizations
Notable scholarly horological organizations include:
American Watchmakers-Clockmakers Institute – AWCI (United States of America)
Antiquarian Horological Society – AHS (United Kingdom)
British Horological Institute – BHI (United Kingdom)
Chronometrophilia (Switzerland)
Deutsche Gesellschaft für Chronometrie – DGC (Germany)
Horological Society of New York – HSNY (United States of America)
National Association of Watch and Clock Collectors – NAWCC (United States of America)
UK Horology - UK Clock & Watch Company based in Bristol
Glossary
See also
Complication (horology)
Hora (astrology)
List of clock manufacturers
List of watch manufacturers
Winthrop Kellogg Edey
Allan variance
Clock drift
International Earth Rotation and Reference Systems Service
Time and Frequency Standards Laboratory
Time deviation
Notes
References
Further reading
Berner, G.A., Illustrated Professional Dictionary of Horology, Federation of the Swiss Watch Industry FH 1961 - 2012
Daniels, George, Watchmaking, London: Philip Wilson Publishers, 1981 (reprinted June 15, 2011)
Beckett, Edmund, A Rudimentary Treatise on Clocks, Watches and Bells, 1903, from Project Gutenberg
Grafton, Edward, Horology, a popular sketch of clock and watch making, London: Aylett and Jones, 1849
Time
Frequency
Metrology
Timekeeping | Chronometry | [
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"Mathematics"
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16,291,643 | https://en.wikipedia.org/wiki/HSL%20%28Fortran%20library%29 | HSL, originally the Harwell Subroutine Library, is a collection of Fortran 77 and 95 codes that address core problems in numerical analysis. It is primarily developed by the Numerical Analysis Group at the Rutherford Appleton Laboratory with contributions from other experts in the field.
HSL codes are easily recognizable by the format of their names, consisting of two letters followed by two numbers, dating back to early versions of Fortran's limited subroutine name length. The letters denote a broad classification of the problem they solve, and the numbers serve to distinguish different codes. For example, the well known sparse LU code MA28 (superseded by MA48) is a Matrix Algebra code number 28. Fortran 95 codes are differentiated from Fortran 77 codes by the prefix HSL_.
History
Early history
Original development of the Harwell Subroutine Library began in 1963 by Mike Powell and Mike Hopper for internal use on an IBM mainframe at AERE Harwell. Early contributors also included Alan Curtis. With a spreading reputation, the Library was distributed externally for the first time in 1964 upon request. The first library catalog (AERE Report M-1748) was released in 1966.
Recent history
Over the intervening years, HSL has striven to maintain a high standard of reliability and has garnered a worldwide reputation as a prime source of numerical software. It has undergone a number of changes to reflect newly available features of the Fortran language, completing in 1990 the conversion to Fortran 77, and more recently, the entire Library has been made thread safe. Many newer codes are written in Fortran 95.
New packages continue to be developed, with a new release issued every two to three years. Many older codes have now been superseded and are available in the HSL Archive.
Licensing
The current version, HSL 2007 is a commercial product sold by AspenTech, but is also available without charge to individual academics direct from STFC for teaching and their own academic research purposes. HSL is currently not sold to commercial competitors of Aspen Technology.
Obsolete routines are stored in the HSL archive and are available for personal non-commercial use by anyone following registration with HSL. Commercial use and distribution of these routines still requires a purchased licence.
References
J.K.Reid and J.A.Scott (Dec 2006, Sep 2007), Guidelines for the development of HSL software, Technical Report RAL-TR-2006-031
M.J.D.Powell 25 years of Theoretical Physics 1954-1979: Chapter XVIII: Numerical Analysis. A special publication by Harwell Research Laboratory of UKAEA
Footnotes
External links
HSL home page at STFC
HSL home page at AspenTech
HSL Archive
Fortran libraries
Numerical software
Science and Technology Facilities Council
Science and technology in Oxfordshire
Vale of White Horse | HSL (Fortran library) | [
"Mathematics"
] | 564 | [
"Numerical software",
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16,297,028 | https://en.wikipedia.org/wiki/GenerativeComponents | GenerativeComponents is parametric CAD software developed by Bentley Systems, was first introduced in 2003, became increasingly used in practice (especially by the London architectural community) by early 2005, and was commercially released in November 2007. GenerativeComponents has a strong traditional base of users in academia and at technologically advanced design firms. GenerativeComponents is often referred to by the nickname of 'GC'. GC epitomizes the quest to bring parametric modeling capabilities of 3D solid modeling into architectural design, seeking to provide greater fluidity and fluency than mechanical 3D solid modeling.
Users can interact with the software by either dynamically modeling and directly manipulating geometry, or by applying rules and capturing relationships among model elements, or by defining complex forms and systems through concisely expressed algorithms.
The software supports many industry standard file input and outputs including DGN by Bentley Systems, DWG by Autodesk, STL (Stereo Lithography), Rhino, and others. The software can also integrate with Building Information Modeling systems, specifically and an installed extension/Companion Feature to Bentley's AECOsim Building Designer.
The software has a published API and uses a simple scripting language, both allowing the integration with many different software tools, and the creation of custom programs by users.
This software is primarily used by architects and engineers in the design of buildings, but has also been used to model natural and biological structures and mathematical systems.
Generative Components currently runs exclusively on Microsoft Windows operating systems, and in English.
Bentley Systems Incorporated offers GC as a free download. This version of GC does not time-out and is not feature limited. It requires registration with an email address. This is a standalone version of GC that includes the underlying Bentley MicroStation software that is required for it to run.
SmartGeometry Group
The SmartGeometry Group has been instrumental in the formation of GenerativeComponents.
GenerativeComponents was brought to the market after utilizing a multi-year testing cycle with a dedicated user community in the SmartGeometry group. This community was responsible for shaping the product very early in its life and continues to play an important role in defining it. The SmartGeometry Group is an independent non-profit organization; it is not a Bentley user group.
The SmartGeometry Group organizes an annual multi-day workshop and accompanying conference highlighting advanced design practices and technology. Recent workshop and conferences have been in:
Munich (2008);
San Francisco (2009);
IAAC - Barcelona, Spain (2010);
CITA - Copenhagen, Denmark (2011);
RPI - Troy, New York (2012);
UCL - London, UK (2013);
Hong Kong (2014);
Gothenburg (2016);
Toronto (2018).
See also
Architecture
Architectural engineering
Design computing
Comparison of CAD Software
References
Computer-aided design
Building engineering software
Computer-aided design software
Data modeling | GenerativeComponents | [
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16,297,790 | https://en.wikipedia.org/wiki/Shewanella%20oneidensis | Shewanella oneidensis is a bacterium notable for its ability to reduce metal ions and live in environments with or without oxygen. This proteobacterium was first isolated from Lake Oneida, NY in 1988, hence its name.
Shewanella oneidensis is a facultative bacterium, capable of surviving and proliferating in both aerobic and anaerobic conditions. The special interest in S. oneidensis MR-1 revolves around its behavior in an anaerobic environment contaminated by heavy metals such as iron, lead and uranium. Experiments suggest it may reduce ionic mercury to elemental mercury and ionic silver to elemental silver. Cellular respiration for these bacteria is not restricted to heavy metals though; the bacteria can also target sulfates, nitrates and chromates when grown anaerobically.
Name
This species is referred to as S. oneidensis MR-1, indicating "manganese reducing", a special feature of this organism. It is a common misconception to think that MR-1 refers to "metal-reducing" instead of the original intended "manganese-reducing" as observed by Kenneth H. Nealson, who first isolated the organism. Although it was originally known as "manganese-reducing", the additional abbreviation expansion of "metal-reducing" is also valid as S. oneidensis MR-1 does reduce metals other than manganese.
Qualities
Metal reduction
Shewanella oneidensis MR-1 belongs to a class of bacteria known as "Dissimilatory Metal-Reducing Bacteria (DMRB)" because of their ability to couple metal reduction with their metabolism. The means of reducing the metals is of particular controversy, as research using scanning electron microscopy and transmission electron microscopy revealed abnormal structural protrusions resembling bacterial filaments that are thought to be involved in the metal reduction. This process of producing an external filament is completely absent from conventional bacterial respiration and is the center of many current studies.
The mechanics of this bacterium's resistance and use of heavy metal ions is deeply related to its metabolism pathway web. Putative multidrug efflux transporters, detoxification proteins, extracytoplasmic sigma factors and PAS domain regulators are shown to have higher expression activity in presence of heavy metal. Cytochrome c class protein SO3300 also has an elevated transcription. For example, when reducing U(VI), special cytochromes such as MtrC and OmcA are used to form UO2 nanoparticles and associate it with biopolymers.
Chemical modification
In 2017 researchers used a synthetic molecule called DSFO+ to modify cell membranes in two mutant strains of Shewanella. DSFO+ could completely replace natural current-conducting proteins, boosting the power that the microbe generated. The process was a chemical modification only that did not modify the organism's genome and that was divided among the bacteria's offspring, diluting the effect.
Pellicle formation
Pellicle is a variety of biofilm that is formed between the air and the liquid in which bacteria grow. In a biofilm, bacterial cells interact with each other to protect their community and co-operate metabolically (microbial communities). In S. oneidensis, pellicle formation is typical and is related to the process of reducing heavy metal. Pellicle formation is extensively researched in this species. Pellicle is usually formed in three steps: cells attach to the triple surface of culture device, air and liquid, then developing a one-layered biofilm from the initial cells, and subsequently maturing to a complicated three-dimensional structure. In a developed pellicle, a number of substances between the cells (extracellular polymeric substances) help maintain the pellicle matrix. The process of pellicle formation involves significant microbial activities and related substances. For the extracellular polymeric substances, many proteins and other bio-macromolecules are required.
Many metal cations are also required in the process. EDTA control and extensive cation presence/absence tests show that Ca(II), Mn(II), Cu(II) and Zn(II) are all essential in this process, probably functioning as a part of a coenzyme or prosthetic group. Mg(II) has partial effect, while Fe(II) and Fe(III) are inhibitory to some degree. Flagella are considered to contribute to pellicle formation. The biofilm needs bacterial cells to move in a certain manner, while flagella is the organelle which has locomotive function. Mutant strains lacking flagella can still form pellicle, albeit much less rapidly.
Applications
Nanotechnology
Shewanella oneidensis MR-1 can change the oxidation state of metals. These microbial processes allow exploration of novel applications, for example, the biosynthesis of metal nanomaterials. In contrast to chemical and physical methods, microbial processes for synthesizing nanomaterials can be achieved in aqueous phase under gentle and environmentally benign conditions. Many organisms can be utilized to synthesize metal nanomaterials. S. oneidensis is able to reduce a diverse range of metal ions extracellularly and this extracellular production greatly facilitates the extraction of nanomaterials. The extracellular electron transport chains responsible for transferring electrons across cell membranes are relatively well characterized, in particular outer membrane c-type cytochromes MtrC and OmcA. A 2013 study suggested that it is possible to alter particle size and activity of extracellular biogenic nanoparticles via controlled expression of the genes encoding surface proteins. An important example is the synthesis of silver nanoparticle by S. oneidensis, where its antibacterial activity can be influenced by the expression of outer membrane c-type cytochromes. Silver nanoparticles are considered to be a new generation of antimicrobial as they exhibit biocidal activity towards a broad range of bacteria, and are gaining importance with the increasing resistance in antibiotics by pathogenic bacteria. Shewanella has been seen in laboratory settings to bioreduce a substantial amount of palladium and dechlorinate near 70% of polychlorinated biphenyls (PCBs). The production of nanoparticles by S. oneidensis MR-1 are closely associated to the MTR pathway (e.g. silver nanoparticles), or the hydrogenase pathway (e.g. palladium nanoparticles).
Wastewater treatment
Shewanella oneidensis' ability to reduce and absorb heavy metals makes it a candidate for use in wastewater treatment.
DSFO+ could possibly allow the bacteria to electrically communicate with an electrode and generate electricity in a wastewater application.
Genome
As a facultative anaerobe with a branching electron transport pathway, S. oneidensis is considered a model organism in microbiology. In 2002, its genomic sequence was published. It has a 4.9Mb circular chromosome that is predicted to encode 4,758 protein open reading frames. It has a 161kb plasmid with 173 open reading frames. A re-annotation was made in 2003.
References
External links
New bacterial behavior observed PNAS study documents puzzling movement of electricity-producing bacteria near energy sources, abstract at Eurekalert
'Rock-Breathing' Bacteria Could Generate Electricity and Clean Up Oil Spills, ScienceDaily (Dec. 15, 2009)
Bacteria that can form electric circuits?
Type strain of Shewanella oneidensis at BacDive – the Bacterial Diversity Metadatabase
Alteromonadales
Pollution control technologies | Shewanella oneidensis | [
"Chemistry",
"Engineering"
] | 1,565 | [
"Pollution control technologies",
"Environmental engineering"
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16,300,571 | https://en.wikipedia.org/wiki/Computational%20creativity | Computational creativity (also known as artificial creativity, mechanical creativity, creative computing or creative computation) is a multidisciplinary endeavour that is located at the intersection of the fields of artificial intelligence, cognitive psychology, philosophy, and the arts (e.g., computational art as part of computational culture).
The goal of computational creativity is to model, simulate or replicate creativity using a computer, to achieve one of several ends:
To construct a program or computer capable of human-level creativity.
To better understand human creativity and to formulate an algorithmic perspective on creative behavior in humans.
To design programs that can enhance human creativity without necessarily being creative themselves.
The field of computational creativity concerns itself with theoretical and practical issues in the study of creativity. Theoretical work on the nature and proper definition of creativity is performed in parallel with practical work on the implementation of systems that exhibit creativity, with one strand of work informing the other.
The applied form of computational creativity is known as media synthesis.
Theoretical issues
Theoretical approaches concern the essence of creativity. Especially, under what circumstances it is possible to call the model a "creative" if eminent creativity is about rule-breaking or the disavowal of convention. This is a variant of Ada Lovelace's objection to machine intelligence, as recapitulated by modern theorists such as Teresa Amabile. If a machine can do only what it was programmed to do, how can its behavior ever be called creative?
Indeed, not all computer theorists would agree with the premise that computers can only do what they are programmed to do—a key point in favor of computational creativity.
Defining creativity in Computational terms
Because no single perspective or definition seems to offer a complete picture of creativity, the AI researchers Newell, Shaw and Simon developed the combination of novelty and usefulness into the cornerstone of a multi-pronged view of creativity, one that uses the following four criteria to categorize a given answer or solution as creative:
The answer is novel and useful (either for the individual or for society)
The answer demands that we reject ideas we had previously accepted
The answer results from intense motivation and persistence
The answer comes from clarifying a problem that was originally vague
Margaret Boden focused on the first two of these criteria, arguing instead that creativity (at least when asking whether computers could be creative) should be defined as "the ability to come up with ideas or artifacts that are new, surprising, and valuable".
Mihali Csikszentmihalyi argued that creativity had to be considered instead in a social context, and his DIFI (Domain-Individual-Field Interaction) framework has since strongly influenced the field. In DIFI, an individual produces works whose novelty and value are assessed by the field—other people in society—providing feedback and ultimately adding the work, now deemed creative, to the domain of societal works from which an individual might be later influenced.
Whereas the above reflects a top-down approach to computational creativity, an alternative thread has developed among bottom-up computational psychologists involved in artificial neural network research. During the late 1980s and early 1990s, for example, such generative neural systems were driven by genetic algorithms. Experiments involving recurrent nets were successful in hybridizing simple musical melodies and predicting listener expectations.
Machine learning for Computational creativity
While traditional computational approaches to creativity rely on the explicit formulation of prescriptions by developers and a certain degree of randomness in computer programs, machine learning methods allow computer programs to learn on heuristics from input data enabling creative capacities within the computer programs. Especially, deep artificial neural networks allow to learn patterns from input data that allow for the non-linear generation of creative artefacts. Before 1989, artificial neural networks have been used to model certain aspects of creativity. Peter Todd (1989) first trained a neural network to reproduce musical melodies from a training set of musical pieces. Then he used a change algorithm to modify the network's input parameters. The network was able to randomly generate new music in a highly uncontrolled manner. In 1992, Todd extended this work, using the so-called distal teacher approach that had been developed by Paul Munro, Paul Werbos, D. Nguyen and Bernard Widrow, Michael I. Jordan and David Rumelhart. In the new approach, there are two neural networks, one of which is supplying training patterns to another.
In later efforts by Todd, a composer would select a set of melodies that define the melody space, position them on a 2-d plane with a mouse-based graphic interface, and train a connectionist network to produce those melodies, and listen to the new "interpolated" melodies that the network generates corresponding to intermediate points in the 2-d plane.
Key concepts from literature
Some high-level and philosophical themes recur throughout the field of computational creativity, for example as follows.
Important categories of creativity
Margaret Boden refers to creativity that is novel merely to the agent that produces it as "P-creativity" (or "psychological creativity"), and refers to creativity that is recognized as novel by society at large as "H-creativity" (or "historical creativity").
Exploratory and transformational creativity
Boden also distinguishes between the creativity that arises from an exploration within an established conceptual space, and the creativity that arises from a deliberate transformation or transcendence of this space. She labels the former as exploratory creativity and the latter as transformational creativity, seeing the latter as a form of creativity far more radical, challenging, and rarer than the former. Following the criteria from Newell and Simon elaborated above, we can see that both forms of creativity should produce results that are appreciably novel and useful (criterion 1), but exploratory creativity is more likely to arise from a thorough and persistent search of a well-understood space (criterion 3) -- while transformational creativity should involve the rejection of some of the constraints that define this space (criterion 2) or some of the assumptions that define the problem itself (criterion 4). Boden's insights have guided work in computational creativity at a very general level, providing more an inspirational touchstone for development work than a technical framework of algorithmic substance. However, Boden's insights are also the subject of formalization, most notably in the work by Geraint Wiggins.
Generation and evaluation
The criterion that creative products should be novel and useful means that creative computational systems are typically structured into two phases, generation and evaluation. In the first phase, novel (to the system itself, thus P-Creative) constructs are generated; unoriginal constructs that are already known to the system are filtered at this stage. This body of potentially creative constructs is then evaluated, to determine which are meaningful and useful and which are not. This two-phase structure conforms to the Geneplore model of Finke, Ward and Smith, which is a psychological model of creative generation based on empirical observation of human creativity.
Co-creation
While much of computational creativity research focuses on independent and automatic machine-based creativity generation, many researchers are inclined towards a collaboration approach. This human-computer interaction is sometimes categorized under the creativity support tools development. These systems aim to provide an ideal framework for research, integration, decision-making, and idea generation. Recently, deep learning approaches to imaging, sound and natural language processing, resulted in the modeling of productive creativity development frameworks.
Innovation
Computational creativity is increasingly being discussed in the innovation and management literature as the recent development in AI may disrupt entire innovation processes and fundamentally change how innovations will be created. Philip Hutchinson highlights the relevance of computational creativity for creating innovation and introduced the concept of “self-innovating artificial intelligence” (SAI) to describe how companies make use of AI in innovation processes to enhance their innovative offerings. SAI is defined as the organizational utilization of AI with the aim of incrementally advancing existing or developing new products, based on insights from continuously combining and analyzing multiple data sources. As AI becomes a general-purpose technology, the spectrum of products to be developed with SAI will broaden from simple to increasingly complex. This implies that computational creativity leads to a shift of creativity-related skills for humans.
Combinatorial creativity
A great deal, perhaps all, of human creativity can be understood as a novel combination of pre-existing ideas or objects. Common strategies for combinatorial creativity include:
Placing a familiar object in an unfamiliar setting (e.g., Marcel Duchamp's Fountain) or an unfamiliar object in a familiar setting (e.g., a fish-out-of-water story such as The Beverly Hillbillies)
Blending two superficially different objects or genres (e.g., a sci-fi story set in the Wild West, with robot cowboys, as in Westworld, or the reverse, as in Firefly; Japanese haiku poems, etc.)
Comparing a familiar object to a superficially unrelated and semantically distant concept (e.g., "Makeup is the Western burka"; "A zoo is a gallery with living exhibits")
Adding a new and unexpected feature to an existing concept (e.g., adding a scalpel to a Swiss Army knife; adding a camera to a mobile phone)
Compressing two incongruous scenarios into the same narrative to get a joke (e.g., the Emo Philips joke "Women are always using men to advance their careers. Damned anthropologists!")
Using an iconic image from one domain in a domain for an unrelated or incongruous idea or product (e.g., using the Marlboro Man image to sell cars, or to advertise the dangers of smoking-related impotence).
The combinatorial perspective allows us to model creativity as a search process through the space of possible combinations. The combinations can arise from composition or concatenation of different representations, or through a rule-based or stochastic transformation of initial and intermediate representations. Genetic algorithms and neural networks can be used to generate blended or crossover representations that capture a combination of different inputs.
Conceptual blending
Mark Turner and Gilles Fauconnier propose a model called Conceptual Integration Networks that elaborates upon Arthur Koestler's ideas about creativity as well as work by Lakoff and Johnson, by synthesizing ideas from Cognitive Linguistic research into mental spaces and conceptual metaphors. Their basic model defines an integration network as four connected spaces:
A first input space (contains one conceptual structure or mental space)
A second input space (to be blended with the first input)
A generic space of stock conventions and image-schemas that allow the input spaces to be understood from an integrated perspective
A blend space in which a selected projection of elements from both input spaces are combined; inferences arising from this combination also reside here, sometimes leading to emergent structures that conflict with the inputs.
Fauconnier and Turner describe a collection of optimality principles that are claimed to guide the construction of a well-formed integration network. In essence, they see blending as a compression mechanism in which two or more input structures are compressed into a single blend structure. This compression operates on the level of conceptual relations. For example, a series of similarity relations between the input spaces can be compressed into a single identity relationship in the blend.
Some computational success has been achieved with the blending model by extending pre-existing computational models of analogical mapping that are compatible by virtue of their emphasis on connected semantic structures. In 2006, Francisco Câmara Pereira presented an implementation of blending theory that employs ideas both from symbolic AI and genetic algorithms to realize some aspects of blending theory in a practical form; his example domains range from the linguistic to the visual, and the latter most notably includes the creation of mythical monsters by combining 3-D graphical models.
Linguistic creativity
Language provides continuous opportunity for creativity, evident in the generation of novel sentences, phrasings, puns, neologisms, rhymes, allusions, sarcasm, irony, similes, metaphors, analogies, witticisms, and jokes. Native speakers of morphologically rich languages frequently create new word-forms that are easily understood, and some have found their way to the dictionary. The area of natural language generation has been well studied, but these creative aspects of everyday language have yet to be incorporated with any robustness or scale.
Hypothesis of creative patterns
In the seminal work of applied linguist Ronald Carter, he hypothesized two main creativity types involving words and word patterns: pattern-reforming creativity, and pattern-forming creativity. Pattern-reforming creativity refers to creativity by the breaking of rules, reforming and reshaping patterns of language often through individual innovation, while pattern-forming creativity refers to creativity via conformity to language rules rather than breaking them, creating convergence, symmetry and greater mutuality between interlocutors through their interactions in the form of repetitions.
Story generation
Substantial work has been conducted in this area of linguistic creation since the 1970s, with the development of James Meehan's TALE-SPIN
system. TALE-SPIN viewed stories as narrative descriptions of a problem-solving effort, and created stories by first establishing a goal for the story's characters so that their search for a solution could be tracked and recorded. The MINSTREL system represents a complex elaboration of this basic approach, distinguishing a range of character-level goals in the story from a range of author-level goals for the story. Systems like Bringsjord's BRUTUS elaborate these ideas further to create stories with complex interpersonal themes like betrayal. Nonetheless, MINSTREL explicitly models the creative process with a set of Transform Recall Adapt Methods (TRAMs) to create novel scenes from old. The MEXICA model of Rafael Pérez y Pérez and Mike Sharples is more explicitly interested in the creative process of storytelling, and implements a version of the engagement-reflection cognitive model of creative writing.
Metaphor and simile
Example of a metaphor: "She was an ape."
Example of a simile: "Felt like a tiger-fur blanket."
The computational study of these phenomena has mainly focused on interpretation as a knowledge-based process. Computationalists such as Yorick Wilks, James Martin, Dan Fass, John Barnden, and Mark Lee have developed knowledge-based approaches to the processing of metaphors, either at a linguistic level or a logical level. Tony Veale and Yanfen Hao have developed a system, called Sardonicus, that acquires a comprehensive database of explicit similes from the web; these similes are then tagged as bona-fide (e.g., "as hard as steel") or ironic (e.g., "as hairy as a bowling ball", "as pleasant as a root canal"); similes of either type can be retrieved on demand for any given adjective. They use these similes as the basis of an on-line metaphor generation system called Aristotle that can suggest lexical metaphors for a given descriptive goal (e.g., to describe a supermodel as skinny, the source terms "pencil", "whip", "whippet", "rope", "stick-insect" and "snake" are suggested).
Analogy
The process of analogical reasoning has been studied from both a mapping and a retrieval perspective, the latter being key to the generation of novel analogies. The dominant school of research, as advanced by Dedre Gentner, views analogy as a structure-preserving process; this view has been implemented in the structure mapping engine or SME, the MAC/FAC retrieval engine (Many Are Called, Few Are Chosen), ACME (Analogical Constraint Mapping Engine) and ARCS (Analogical Retrieval Constraint System). Other mapping-based approaches include Sapper, which situates the mapping process in a semantic-network model of memory. Analogy is a very active sub-area of creative computation and creative cognition; active figures in this sub-area include Douglas Hofstadter, Paul Thagard, and Keith Holyoak. Also worthy of note here is Peter Turney and Michael Littman's machine learning approach to the solving of SAT-style analogy problems; their approach achieves a score that compares well with average scores achieved by humans on these tests.
Joke generation
Humour is an especially knowledge-hungry process, and the most successful joke-generation systems to date have focussed on pun-generation, as exemplified by the work of Kim Binsted and Graeme Ritchie. This work includes the JAPE system, which can generate a wide range of puns that are consistently evaluated as novel and humorous by young children. An improved version of JAPE has been developed in the guise of the STANDUP system, which has been experimentally deployed as a means of enhancing linguistic interaction with children with communication disabilities. Some limited progress has been made in generating humour that involves other aspects of natural language, such as the deliberate misunderstanding of pronominal reference (in the work of Hans Wim Tinholt and Anton Nijholt), as well as in the generation of humorous acronyms in the HAHAcronym system of Oliviero Stock and Carlo Strapparava.
Neologism
The blending of multiple word forms is a dominant force for new word creation in language; these new words are commonly called "blends" or "portmanteau words" (after Lewis Carroll). Tony Veale has developed a system called ZeitGeist that harvests neological headwords from Wikipedia and interprets them relative to their local context in Wikipedia and relative to specific word senses in WordNet. ZeitGeist has been extended to generate neologisms of its own; the approach combines elements from an inventory of word parts that are harvested from WordNet, and simultaneously determines likely glosses for these new words (e.g., "food traveller" for "gastronaut" and "time traveller" for "chrononaut"). It then uses Web search to determine which glosses are meaningful and which neologisms have not been used before; this search identifies the subset of generated words that are both novel ("H-creative") and useful.
A corpus linguistic approach to the search and extraction of neologism have also shown to be possible. Using Corpus of Contemporary American English as a reference corpus, Locky Law has performed an extraction of neologism, portmanteaus and slang words using the hapax legomena which appeared in the scripts of American TV drama House M.D.
In terms of linguistic research in neologism, Stefan Th. Gries has performed a quantitative analysis of blend structure in English and found that "the degree of recognizability of the source words and that the similarity of source words to the blend plays a vital role in blend formation." The results were validated through a comparison of intentional blends to speech-error blends.
Poetry
Like jokes, poems involve a complex interaction of different constraints, and no general-purpose poem generator adequately combines the meaning, phrasing, structure and rhyme aspects of poetry. Nonetheless, Pablo Gervás has developed a noteworthy system called ASPERA that employs a case-based reasoning (CBR) approach to generating poetic formulations of a given input text via a composition of poetic fragments that are retrieved from a case-base of existing poems. Each poem fragment in the ASPERA case-base is annotated with a prose string that expresses the meaning of the fragment, and this prose string is used as the retrieval key for each fragment. Metrical rules are then used to combine these fragments into a well-formed poetic structure. Racter is an example of such a software project.
Musical creativity
Computational creativity in the music domain has focused both on the generation of musical scores for use by human musicians, and on the generation of music for performance by computers. The domain of generation has included classical music (with software that generates music in the style of Mozart and Bach) and jazz. Most notably, David Cope has written a software system called "Experiments in Musical Intelligence" (or "EMI") that is capable of analyzing and generalizing from existing music by a human composer to generate novel musical compositions in the same style. EMI's output is convincing enough to persuade human listeners that its music is human-generated to a high level of competence.
In the field of contemporary classical music, Iamus is the first computer that composes from scratch, and produces final scores that professional interpreters can play. The London Symphony Orchestra played a piece for full orchestra, included in Iamus' debut CD, which New Scientist described as "The first major work composed by a computer and performed by a full orchestra". Melomics, the technology behind Iamus, is able to generate pieces in different styles of music with a similar level of quality.
Creativity research in jazz has focused on the process of improvisation and the cognitive demands that this places on a musical agent: reasoning about time, remembering and conceptualizing what has already been played, and planning ahead for what might be played next.
The robot Shimon, developed by Gil Weinberg of Georgia Tech, has demonstrated jazz improvisation. Virtual improvisation software based on researches on stylistic modeling carried out by Gerard Assayag and Shlomo Dubnov include OMax, SoMax and PyOracle, are used to create improvisations in real-time by re-injecting variable length sequences learned on the fly from the live performer.
In the field of musical composition, the patented works by René-Louis Baron allowed to make a robot that can create and play a multitude of orchestrated melodies, so-called "coherent" in any musical style. All outdoor physical parameter associated with one or more specific musical parameters, can influence and develop each of these songs (in real-time while listening to the song). The patented invention Medal-Composer raises problems of copyright.
Visual and artistic creativity
Computational creativity in the generation of visual art has had some notable successes in the creation of both abstract art and representational art. A well-known program in this domain is Harold Cohen's AARON, which has been continuously developed and augmented since 1973. Though formulaic, Aaron exhibits a range of outputs, generating black-and-white drawings or colour paintings that incorporate human figures (such as dancers), potted plants, rocks, and other elements of background imagery. These images are of a sufficiently high quality to be displayed in reputable galleries.
Other software artists of note include the NEvAr system (for "Neuro-Evolutionary Art") of Penousal Machado. NEvAr uses a genetic algorithm to derive a mathematical function that is then used to generate a coloured three-dimensional surface. A human user is allowed to select the best pictures after each phase of the genetic algorithm, and these preferences are used to guide successive phases, thereby pushing NEvAr's search into pockets of the search space that are considered most appealing to the user.
The Painting Fool, developed by Simon Colton originated as a system for overpainting digital images of a given scene in a choice of different painting styles, colour palettes and brush types. Given its dependence on an input source image to work with, the earliest iterations of the Painting Fool raised questions about the extent of, or lack of, creativity in a computational art system. Nonetheless, The Painting Fool has been extended to create novel images, much as AARON does, from its own limited imagination. Images in this vein include cityscapes and forests, which are generated by a process of constraint satisfaction from some basic scenarios provided by the user (e.g., these scenarios allow the system to infer that objects closer to the viewing plane should be larger and more color-saturated, while those further away should be less saturated and appear smaller). Artistically, the images now created by the Painting Fool appear on a par with those created by Aaron, though the extensible mechanisms employed by the former (constraint satisfaction, etc.) may well allow it to develop into a more elaborate and sophisticated painter.
The artist Krasi Dimtch (Krasimira Dimtchevska) and the software developer Svillen Ranev have created a computational system combining a rule-based generator of English sentences and a visual composition builder that converts sentences generated by the system into abstract art. The software generates automatically indefinite number of different images using different color, shape and size palettes. The software also allows the user to select the subject of the generated sentences or/and the one or more of the palettes used by the visual composition builder.
An emerging area of computational creativity is that of video games. ANGELINA is a system for creatively developing video games in Java by Michael Cook. One important aspect is Mechanic Miner, a system that can generate short segments of code that act as simple game mechanics. ANGELINA can evaluate these mechanics for usefulness by playing simple unsolvable game levels and testing to see if the new mechanic makes the level solvable. Sometimes Mechanic Miner discovers bugs in the code and exploits these to make new mechanics for the player to solve problems with.
In July 2015, Google released DeepDream – an open source computer vision program, created to detect faces and other patterns in images with the aim of automatically classifying images, which uses a convolutional neural network to find and enhance patterns in images via algorithmic pareidolia, thus creating a dreamlike psychedelic appearance in the deliberately over-processed images.
In August 2015, researchers from Tübingen, Germany created a convolutional neural network that uses neural representations to separate and recombine content and style of arbitrary images which is able to turn images into stylistic imitations of works of art by artists such as a Picasso or Van Gogh in about an hour. Their algorithm is put into use in the website DeepArt that allows users to create unique artistic images by their algorithm.
In early 2016, a global team of researchers explained how a new computational creativity approach known as the Digital Synaptic Neural Substrate (DSNS) could be used to generate original chess puzzles that were not derived from endgame databases. The DSNS is able to combine features of different objects (e.g. chess problems, paintings, music) using stochastic methods in order to derive new feature specifications which can be used to generate objects in any of the original domains. The generated chess puzzles have also been featured on YouTube.
Creativity in problem solving
Creativity is also useful in allowing for unusual solutions in problem solving. In psychology and cognitive science, this research area is called creative problem solving. The Explicit-Implicit Interaction (EII) theory of creativity has been implemented using a CLARION-based computational model that allows for the simulation of incubation and insight in problem-solving. The emphasis of this computational creativity project is not on performance per se (as in artificial intelligence projects) but rather on the explanation of the psychological processes leading to human creativity and the reproduction of data collected in psychology experiments. So far, this project has been successful in providing an explanation for incubation effects in simple memory experiments, insight in problem solving, and reproducing the overshadowing effect in problem solving.
Debate about "general" theories of creativity
Some researchers feel that creativity is a complex phenomenon whose study is further complicated by the plasticity of the language we use to describe it. We can describe not just the agent of creativity as "creative" but also the product and the method. Consequently, it could be claimed that it is unrealistic to speak of a general theory of creativity. Nonetheless, some generative principles are more general than others, leading some advocates to claim that certain computational approaches are "general theories". Stephen Thaler, for instance, proposes that certain modalities of neural networks are generative enough, and general enough, to manifest a high degree of creative capabilities.
Criticism of computational creativity
Traditional computers, as mainly used in the computational creativity application, do not support creativity, as they fundamentally transform a set of discrete, limited domain of input parameters into a set of discrete, limited domain of output parameters using a limited set of computational functions. As such, a computer cannot be creative, as everything in the output must have been already present in the input data or the algorithms. Related discussions and references to related work are captured in work on philosophical foundations of simulation.
Mathematically, the same set of arguments against creativity has been made by Chaitin. Similar observations come from a Model Theory perspective. All this criticism emphasizes that computational creativity is useful and may look like creativity, but it is not real creativity, as nothing new is created, just transformed in well-defined algorithms.
Events
The International Conference on Computational Creativity (ICCC) occurs annually, organized by The Association for Computational Creativity. Events in the series include:
ICCC 2023: University of Waterloo in Ontario, Canada
ICCC 2022: Free University of Bozen-Bolzano, Bolzano, Italy
ICCC 2021: Mexico City, Mexico (Virtual due to COVID-19 pandemic)
ICCC 2020, Coimbra, Portugal (Virtual due to COVID-19 pandemic)
ICCC 2019, Charlotte, North Carolina, US
ICCC 2018, Salamanca, Spain
ICCC 2017, Atlanta, Georgia, US
ICCC 2016, Paris, France
ICCC 2015, Park City, Utah, US. Keynote: Emily Short
ICCC 2014, Ljubljana, Slovenia. Keynote: Oliver Deussen
ICCC 2013, Sydney, Australia. Keynote: Arne Dietrich
ICCC 2012, Dublin, Ireland. Keynote: Steven Smith
ICCC 2011, Mexico City, Mexico. Keynote: George E Lewis
ICCC 2010, Lisbon, Portugal. Keynote/Invited Talks: Nancy J Nersessian and Mary Lou Maher
Previously, the community of computational creativity has held a dedicated workshop, the International Joint Workshop on Computational Creativity, every year since 1999. Previous events in this series include:
IJWCC 2003, Acapulco, Mexico, as part of IJCAI'2003
IJWCC 2004, Madrid, Spain, as part of ECCBR'2004
IJWCC 2005, Edinburgh, UK, as part of IJCAI'2005
IJWCC 2006, Riva del Garda, Italy, as part of ECAI'2006
IJWCC 2007, London, UK, a stand-alone event
IJWCC 2008, Madrid, Spain, a stand-alone event
The 1st Conference on Computer Simulation of Musical Creativity will be held
CCSMC 2016, 17–19 June, University of Huddersfield, UK. Keynotes: Geraint Wiggins and Graeme Bailey.
See also
1 the Road (1st novel)
Artificial imagination
Algorithmic art
Algorithmic composition
Applications of artificial intelligence
Computer art
Creative computing
Digital morphogenesis
Digital poetry
Generative art
Generative systems
Intrinsic motivation (artificial intelligence)
Musikalisches Würfelspiel (Musical dice game)
Procedural generation
Lists
List of emerging technologies
Outline of artificial intelligence
References
Further reading
An Overview of Artificial Creativity on Think Artificial
Cohen, H., "the further exploits of AARON, Painter" , SEHR, volume 4, issue 2: Constructions of the Mind, 1995
External links
Documentaries
Noorderlicht: Margaret Boden and Stephen Thaler on Creative Computers on Archive.org
In Its Image on Archive.org
Cognitive psychology
Computational fields of study
Creativity techniques
Philosophy of artificial intelligence | Computational creativity | [
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16,304,394 | https://en.wikipedia.org/wiki/DoD%20IPv6%20product%20certification | The United States Department of Defense (DoD) Internet Protocol version 6 (IPv6) product certification program began as a mandate from the DoD's Assistant Secretary of Defense for Networks & Information Integration (ASD-NII) in 2005. The program mandates the Joint Interoperability Test Command (JITC) in Fort Huachuca, Arizona, to test and certify IT products for IPv6 capability according to the RFCs outlined in the DoD's IPv6 Standards Profiles for IPv6 Capable Products. Once products are certified for special interoperability, they are added to the DoD's Unified Capabilities Approved Products List (UC APL) for IPv6. This list is used by procurement offices in the DoD and the U.S. Federal agencies for ongoing purchases and acquisitions of IT equipment.
As of February 2009, the DoD ceased the requirement for IPv6-only testing for certification and entry into the Unified Capabilities Approved Products List (UC APL). According to Kris Strance, DoD CIO IPv6 Lead, "The testing of IPv6 is a part of all product evaluations — it is much broader in scope now." The UC APL is now a single consolidated list of products that have completed Interoperability (IO) and Information Assurance (IA) certification.
DoD IPv6 standards
The DoD IPv6 Standards Profiles for IPv6 Capable Products (DoD IPv6 Profile) is the singular “IPv6 Capable” definition in DoD. It is a document that lists the six agreed upon product classes (Host, Router, Layer 3 Switch, Network Appliance, Security Device, and Advanced Server) and their corresponding standards (RFCs). It lists each standard according to its level of requirement:
MUST: The standard is required to be implemented in the product now.
SHOULD: The standard is optional, but recommended for implementation.
SHOULD+: The standard is optional now, but will be required within a short period of time.
DoD IPv6 generic test plan
The JITC uses its publicly available IPv6 Generic Test Plan (GTP) to test each product for its conformance, performance and interoperability of IPv6 according to the DoD IPv6 Profile. The JITC uses a combination of automated testing tools and manual functional test procedures to conduct this testing.
Process
The vendor, or Program Manager, must make their intentions known to test by providing the JITC with a Letter of Compliance (LoC). This letter will consist of the product to test, the product class it belongs to, a listing of all of the standards that it implements, and a signature from a Vice President or officer of the company. This is the “gateway” to the testing process.
Once the LoC is received, the product is then scheduled for test.
Approximately 6 weeks before the start of testing, the vendor must provide the JITC with funding. This funding must be in the form of a check. The amount is only to charge direct labor hours for testing by the contractor labor support.
If the product successfully meets the criteria, it will be entered on the DoD's UC APL for IPv6.
An IPv6 Capable Special Interoperability Certification Letter and Report will accompany the entry within 30–60 days after testing.
IPv6 pre-certification testing advocates
There are many companies and organizations that help develop and test products for vendors prior to testing at the JITC. These organizations cannot grant certification, but can conduct pre-testing to ensure a vendor's product will pass the necessary certification. Below is a list of these organizations:
The University of New Hampshire InterOperability Laboratory - IPv6 Ready Logo testing:
The IPv6 Ready Logo program
The IPv6 Forum has a service called IPv6 Ready Logo. This service represents a qualification program that assures devices have been tested and are IPv6 capable. Once certified, the service grants qualified products to display their logo. In the IPv6 Forum, they present objectives that are to:
Verify protocol implementation and validate interoperability of IPv6 products.
Provide access to free self-testing tools.
Provide IPv6 Ready Logo testing laboratories across the globe whom will be dedicated to providing testing assistance or services.
IPv6 experts suggest only pursuing to purchase devices given the Phase-2 approval or gold logo since they are given the full treatment:
The Department of Defense (DoD) is committed to IPv6 and will likely be the first federal organization completely converted to IPv6. They also have a process for qualifying IPv6 equipment.
JITC/DISA
The task of certifying IPv6 products was given to the Joint Interoperability Test Command (JITC), part of the Defense Information Systems Agency (DISA). To help standardize IPv6 qualification procedures, the JITC follows what’s called the IPv6 Generic Test Plan.
After JITC qualifies a product, it is added to the Unified Capabilities Approved Products List. Fortunately, JITC makes the list available to the public.
References
The Unified Capabilities Approved Products List for IPv6: https://web.archive.org/web/20081007010445/http://jitc.fhu.disa.mil/apl/ipv6.html
Official IPv6 Capable Certification Testing Process: http://jitc.fhu.disa.mil/apl/ipv6/pdf/ipv6_certification_process_ipv6_v2.pdf
The DoD IPv6 Standards Profiles for IPv6 Capable Products, Version 4: https://web.archive.org/web/20100705061401/http://jitc.fhu.disa.mil/apl/ipv6/pdf/disr_ipv6_product_profile_v4.pdf
The DoD IPv6 Generic Test Plan: https://web.archive.org/web/20100704234338/http://jitc.fhu.disa.mil/adv_ip/register/docs/ipv6v4_may09.pdf
The Testing Times, March 2008, Volume 15, Number 1: http://jitc.fhu.disa.mil/tst_time/docs/year/mar08.pdf
http://www.techrepublic.com/blog/networking/ipv6-capable-devices-make-sure-they-are-ready/2522
Military in Arizona
Interoperability
IPv6 | DoD IPv6 product certification | [
"Engineering"
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16,305,705 | https://en.wikipedia.org/wiki/Milnor%E2%80%93Thurston%20kneading%20theory | The Milnor–Thurston kneading theory is a mathematical theory which analyzes the iterates of piecewise monotone mappings of an interval into itself. The emphasis is on understanding the properties of the mapping that are invariant under topological conjugacy.
The theory had been developed by John Milnor and William Thurston in two widely circulated and influential Princeton preprints from 1977 that were revised in 1981 and finally published in 1988. Applications of the theory include piecewise linear models, counting of fixed points, computing the total variation, and constructing an invariant measure with maximal entropy.
Short description
Kneading theory provides an effective calculus for describing the qualitative behavior of the iterates of a piecewise monotone mapping f of a closed interval I of the real line into itself. Some quantitative invariants of this discrete dynamical system, such as the lap numbers of the iterates and the Artin–Mazur zeta function of f are expressed in terms of certain matrices and formal power series.
The basic invariant of f is its kneading matrix, a rectangular matrix with coefficients in the ring of integer formal power series. A closely related kneading determinant is a formal power series
with odd integer coefficients. In the simplest case when the map is unimodal, with a maximum at c, each coefficient is either or , according to whether the th iterate has local maximum or local minimum at c.
See also
Sharkovsky theorem
Topological entropy
References
Topological dynamics | Milnor–Thurston kneading theory | [
"Mathematics"
] | 305 | [
"Topology",
"Topological dynamics",
"Dynamical systems"
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58,648 | https://en.wikipedia.org/wiki/Faraday%20constant | In physical chemistry, the Faraday constant (symbol , sometimes stylized as ℱ) is a physical constant defined as the quotient of the total electric charge (q) by the amount (n) of elementary charge carriers in any given sample of matter: it is expressed in units of coulombs per mole (C/mol).
As such, it represents the "molar elementary charge", that is, the electric charge of one mole of elementary carriers (e.g., protons). It is named after the English scientist Michael Faraday. Since the 2019 revision of the SI, the Faraday constant has an exactly defined value, the product of the elementary charge (e, in coulombs) and the Avogadro constant (NA, in reciprocal moles):
Derivation
The Faraday constant can be thought of as the conversion factor between the mole (used in chemistry) and the coulomb (used in physics and in practical electrical measurements), and is therefore of particular use in electrochemistry. Because there are exactly NA = entities per mole, and there are exactly elementary charges per coulomb, the Faraday constant is given by the quotient of these two quantities:
One common use of the Faraday constant is in electrolysis calculations. One can divide the amount of charge (the current integrated over time) by the Faraday constant in order to find the chemical amount of a substance (in moles) that has been electrolyzed.
The value of was first determined in the 1800s by weighing the amount of silver deposited in an electrochemical reaction, in which a measured current was passed for a measured time, and using Faraday's law of electrolysis. Until about 1970, the most reliable value of the Faraday constant was determined by a related method of electro-dissolving silver metal in perchloric acid.
Other common units
96.485 kJ per volt–gram-equivalent
23.061 kcal per volt–gram-equivalent
26.801 A·h/mol
Faraday – a unit of charge
Related to the Faraday constant is the "faraday", a unit of electrical charge. Its use is much less common than of the coulomb, but is sometimes used in electrochemistry. One faraday of charge is the charge of one mole of elementary charges (or of negative one mole of electrons), that is,
1 faraday = F × 1 mol = .
Conversely, the Faraday constant F equals 1 faraday per mole.
The faraday is not to be confused with the farad, an unrelated unit of capacitance ().
See also
Farad, the unit of electrical capacitance
Faraday efficiency
Faraday's laws of electrolysis
Faraday cup
References
Electrochemical concepts
Physical constants
Michael Faraday
Units of electrical charge
Units of amount of substance
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58,650 | https://en.wikipedia.org/wiki/Base%20unit%20of%20measurement | A base unit of measurement (also referred to as a base unit or fundamental unit) is a unit of measurement adopted for a base quantity. A base quantity is one of a conventionally chosen subset of physical quantities, where no quantity in the subset can be expressed in terms of the others. The SI base units, or Systéme International d'unités, consists of the metre, kilogram, second, ampere, kelvin, mole and candela.
A unit multiple (or multiple of a unit) is an integer multiple of a given unit; likewise a unit submultiple (or submultiple of a unit) is a submultiple or a unit fraction of a given unit.
Unit prefixes are common base-10 or base-2 powers multiples and submultiples of units.
While a base unit is one that has been explicitly so designated,
a derived unit is unit for a derived quantity, involving the combination of quantities with different units; several SI derived units are specially named.
A coherent derived unit involves no conversion factors.
Background
In the language of measurement, physical quantities are quantifiable aspects of the world, such as time, distance, velocity, mass, temperature, energy, and weight, and units are used to describe their magnitude or quantity. Many of these quantities are related to each other by various physical laws, and as a result the units of a quantities can be generally be expressed as a product of powers of other units; for example, momentum is mass multiplied by velocity, while velocity is distance divided by time. These relationships are discussed in dimensional analysis. Those that can be expressed in this fashion in terms of the base units are called derived units.
International System of Units
In the International System of Units (SI), there are seven base units: kilogram, metre, candela, second, ampere, kelvin, and mole.
Several derived units have been defined, many with special names and symbols.
In 2019 the seven SI base units were redefined in terms of seven defining constants. Therefore the SI base units are no longer necessary but were retained because for historical and practical reasons. See 2019 revision of the SI.
Natural units
A set of base dimensions of quantity is a minimal set of units such that every physical quantity can be expressed in terms of this set. The traditional base dimensions are mass, length, time, charge, and temperature, but in principle, other base quantities could be used. Electric current could be used instead of charge or speed could be used instead of length. Some physicists have not recognized temperature as a base dimension since it simply expresses the energy per particle per degree of freedom which can be expressed in terms of energy (or mass, length, and time). Duff argues that only dimensionless values have physical meaning and all dimensional units are human constructs.
There are other relationships between physical quantities that can be expressed by means of fundamental constants, and to some extent it is an arbitrary decision whether to retain the fundamental constant as a quantity with dimensions or simply to define it as unity or a fixed dimensionless number, and reduce the number of explicit base quantities by one. The ontological issue is whether these fundamental constants really exist as dimensional or dimensionless quantities. This is equivalent to treating length as the same as time or understanding electric charge as a combination of quantities of mass, length, and time which may seem less natural than thinking of temperature as measuring the same material as energy (which is expressible in terms of mass, length, and time).
For instance, time and distance are related to each other by the speed of light, c, which is a fundamental constant. It is possible to use this relationship to eliminate either the base unit of time or that of distance. Similar considerations apply to the Planck constant, h, which relates energy (with dimension expressible in terms of mass, length and time) to frequency (with dimension expressible in terms of time). In theoretical physics it is customary to use such units (natural units) in which and . A similar choice can be applied to the vacuum permittivity, ε0.
One could eliminate either the metre or the second by setting c to unity (or to any other fixed dimensionless number).
One could then eliminate the kilogram by setting ħ to a dimensionless number.
One could eliminate the ampere by setting either the vacuum permittivity ε0 or the elementary charge e to a dimensionless number.
One could eliminate the mole as a base unit by setting the Avogadro constant N to 1. This is natural as it is a technical scaling constant.
One could eliminate the kelvin as it can be argued that temperature simply expresses the energy per particle per degree of freedom, which can be expressed in terms of energy (or mass, length, and time). Another way of saying this is that the Boltzmann constant kB is a technical scaling constant and could be set to a fixed dimensionless number.
Similarly, one could eliminate the candela, as that is defined in terms of other physical quantities via a technical scaling constant, K.
That leaves one base dimension and an associated base unit, but there are several fundamental constants left to eliminate that too – for instance, one could use G, the gravitational constant, me, the electron rest mass, or Λ, the cosmological constant.
The preferred choices vary by the field in physics. Using natural units leaves every physical quantity expressed as a dimensionless number, which is noted by physicists disputing the existence of incompatible base quantities.
See also
Characteristic units
Dimensional analysis
Natural units
One (unit)
References
Measurement
Dimensional analysis
ro:Mărimi fizice fundamentale | Base unit of measurement | [
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58,664 | https://en.wikipedia.org/wiki/Gas%20turbine | A gas turbine or gas turbine engine is a type of continuous flow internal combustion engine. The main parts common to all gas turbine engines form the power-producing part (known as the gas generator or core) and are, in the direction of flow:
a rotating gas compressor
a combustor
a compressor-driving turbine.
Additional components have to be added to the gas generator to suit its application. Common to all is an air inlet but with different configurations to suit the requirements of marine use, land use or flight at speeds varying from stationary to supersonic. A propelling nozzle is added to produce thrust for flight. An extra turbine is added to drive a propeller (turboprop) or ducted fan (turbofan) to reduce fuel consumption (by increasing propulsive efficiency) at subsonic flight speeds. An extra turbine is also required to drive a helicopter rotor or land-vehicle transmission (turboshaft), marine propeller or electrical generator (power turbine). Greater thrust-to-weight ratio for flight is achieved with the addition of an afterburner.
The basic operation of the gas turbine is a Brayton cycle with air as the working fluid: atmospheric air flows through the compressor that brings it to higher pressure; energy is then added by spraying fuel into the air and igniting it so that the combustion generates a high-temperature flow; this high-temperature pressurized gas enters a turbine, producing a shaft work output in the process, used to drive the compressor; the unused energy comes out in the exhaust gases that can be repurposed for external work, such as directly producing thrust in a turbojet engine, or rotating a second, independent turbine (known as a power turbine) that can be connected to a fan, propeller, or electrical generator. The purpose of the gas turbine determines the design so that the most desirable split of energy between the thrust and the shaft work is achieved. The fourth step of the Brayton cycle (cooling of the working fluid) is omitted, as gas turbines are open systems that do not reuse the same air.
Gas turbines are used to power aircraft, trains, ships, electrical generators, pumps, gas compressors, and tanks.
Timeline of development
50: Earliest records of Hero's engine (aeolipile). It most likely served no practical purpose, and was rather more of a curiosity; nonetheless, it demonstrated an important principle of physics that all modern turbine engines rely on.
1000: The "Trotting Horse Lamp" (, zŏumădēng) was used by the Chinese at lantern fairs as early as the Northern Song dynasty. When the lamp is lit, the heated airflow rises and drives an impeller with horse-riding figures attached on it, whose shadows are then projected onto the outer screen of the lantern.
1500: The Smoke jack was drawn by Leonardo da Vinci: Hot air from a fire rises through a single-stage axial turbine rotor mounted in the exhaust duct of the fireplace and turns the roasting spit by gear-chain connection.
1791: A patent was given to John Barber, an Englishman, for the first true gas turbine. His invention had most of the elements present in the modern day gas turbines. The turbine was designed to power a horseless carriage.
1894: Sir Charles Parsons patented the idea of propelling a ship with a steam turbine, and built a demonstration vessel, the Turbinia, easily the fastest vessel afloat at the time.
1899: Charles Gordon Curtis patented the first gas turbine engine in the US.
1900: Sanford Alexander Moss submitted a thesis on gas turbines. In 1903, Moss became an engineer for General Electric's Steam Turbine Department in Lynn, Massachusetts. While there, he applied some of his concepts in the development of the turbocharger.
1903: A Norwegian, Ægidius Elling, built the first gas turbine that was able to produce more power than needed to run its own components, which was considered an achievement in a time when knowledge about aerodynamics was limited. Using rotary compressors and turbines it produced .
1904: A gas turbine engine designed by Franz Stolze, based on his earlier 1873 patent application, is built and tested in Berlin. The Stolze gas turbine was too inefficient to sustain its own operation.
1906: The Armengaud-Lemale gas turbine tested in France. This was a relatively large machine which included a 25-stage centrifugal compressor designed by Auguste Rateau and built by the Brown Boveri Company. The gas turbine could sustain its own air compression but was too inefficient to produce useful work.
1910: The first operational Holzwarth gas turbine (pulse combustion) achieves an output of . Planned output of the machine was and its efficiency is below that of contemporary reciprocating engines.
1920s The practical theory of gas flow through passages was developed into the more formal (and applicable to turbines) theory of gas flow past airfoils by A. A. Griffith resulting in the publishing in 1926 of An Aerodynamic Theory of Turbine Design. Working testbed designs of axial turbines suitable for driving a propeller were developed by the Royal Aeronautical Establishment.
1930: Having found no interest from the RAF for his idea, Frank Whittle patented the design for a centrifugal gas turbine for jet propulsion. The first successful test run of his engine occurred in England in April 1937.
1932: The Brown Boveri Company of Switzerland starts selling axial compressor and turbine turbosets as part of the turbocharged steam generating Velox boiler. Following the gas turbine principle, the steam evaporation tubes are arranged within the gas turbine combustion chamber; the first Velox plant is erected at a French Steel mill in Mondeville, Calvados.
1936: The first constant flow industrial gas turbine is commissioned by the Brown Boveri Company and goes into service at Sun Oil's Marcus Hook refinery in Pennsylvania, US.
1937: Working proof-of-concept prototype turbojet engine runs in UK (Frank Whittle's) and Germany (Hans von Ohain's Heinkel HeS 1). Henry Tizard secures UK government funding for further development of Power Jets engine.
1939: The First 4 MW utility power generation gas turbine is built by the Brown Boveri Company for an emergency power station in Neuchâtel, Switzerland. The turbojet powered Heinkel He 178, the world's first jet aircraft, makes its first flight.
1940: Jendrassik Cs-1, a turboprop engine, made its first bench run. The Cs-1 was designed by Hungarian engineer György Jendrassik, and was intended to power a Hungarian twin-engine heavy fighter, the RMI-1. Work on the Cs-1 stopped in 1941 without the type having powered any aircraft.
1944: The Junkers Jumo 004 engine enters full production, powering the first German military jets such as the Messerschmitt Me 262. This marks the beginning of the reign of gas turbines in the sky.
1946: National Gas Turbine Establishment formed from Power Jets and the RAE turbine division to bring together Whittle and Hayne Constant's work. In Beznau, Switzerland the first commercial reheated/recuperated unit generating 27 MW was commissioned.
1947: A Metropolitan Vickers G1 (Gatric) becomes the first marine gas turbine when it completes sea trials on the Royal Navy's M.G.B 2009 vessel. The Gatric was an aeroderivative gas turbine based on the Metropolitan Vickers F2 jet engine.
1995: Siemens becomes the first manufacturer of large electricity producing gas turbines to incorporate single crystal turbine blade technology into their production models, allowing higher operating temperatures and greater efficiency.
2011 Mitsubishi Heavy Industries tests the first >60% efficiency combined cycle gas turbine (the M501J) at its Takasago, Hyōgo, works.
Theory of operation
In an ideal gas turbine, gases undergo four thermodynamic processes: an isentropic compression, an isobaric (constant pressure) combustion, an isentropic expansion and isobaric heat rejection. Together, these make up the Brayton cycle, also known as the "constant pressure cycle". It is distinguished from the Otto cycle, in that all the processes (compression, ignition combustion, exhaust), occur at the same time, continuously.
In a real gas turbine, mechanical energy is changed irreversibly (due to internal friction and turbulence) into pressure and thermal energy when the gas is compressed (in either a centrifugal or axial compressor). Heat is added in the combustion chamber and the specific volume of the gas increases, accompanied by a slight loss in pressure. During expansion through the stator and rotor passages in the turbine, irreversible energy transformation once again occurs. Fresh air is taken in, in place of the heat rejection.
Air is taken in by a compressor, called a gas generator, with either an axial or centrifugal design, or a combination of the two. This air is then ducted into the combustor section which can be of a annular, can, or can-annular design. In the combustor section, roughly 70% of the air from the compressor is ducted around the combustor itself for cooling purposes. The remaining roughly 30% the air is mixed with fuel and ignited by the already burning air-fuel mixture, which then expands producing power across the turbine. This expansion of the mixture then leaves the combustor section and has its velocity increased across the turbine section to strike the turbine blades, spinning the disc they are attached to, thus creating useful power. Of the power produced, 60-70% is solely used to power the gas generator. The remaining power is used to power what the engine is being used for, typically an aviation application, being thrust in a turbojet, driving the fan of a turbofan, rotor or accessory of a turboshaft, and gear reduction and propeller of a turboprop.
If the engine has a power turbine added to drive an industrial generator or a helicopter rotor, the exit pressure will be as close to the entry pressure as possible with only enough energy left to overcome the pressure losses in the exhaust ducting and expel the exhaust. For a turboprop engine there will be a particular balance between propeller power and jet thrust which gives the most economical operation. In a turbojet engine only enough pressure and energy is extracted from the flow to drive the compressor and other components. The remaining high-pressure gases are accelerated through a nozzle to provide a jet to propel an aircraft.
The smaller the engine, the higher the rotation rate of the shaft must be to attain the required blade tip speed. Blade-tip speed determines the maximum pressure ratios that can be obtained by the turbine and the compressor. This, in turn, limits the maximum power and efficiency that can be obtained by the engine. In order for tip speed to remain constant, if the diameter of a rotor is reduced by half, the rotational speed must double. For example, large jet engines operate around 10,000–25,000 rpm, while micro turbines spin as fast as 500,000 rpm.
Mechanically, gas turbines can be considerably less complex than Reciprocating engines. Simple turbines might have one main moving part, the compressor/shaft/turbine rotor assembly, with other moving parts in the fuel system. This, in turn, can translate into price. For instance, costing for materials, the Jumo 004 proved cheaper than the Junkers 213 piston engine, which was , and needed only 375 hours of lower-skill labor to complete (including manufacture, assembly, and shipping), compared to 1,400 for the BMW 801. This, however, also translated into poor efficiency and reliability. More advanced gas turbines (such as those found in modern jet engines or combined cycle power plants) may have 2 or 3 shafts (spools), hundreds of compressor and turbine blades, movable stator blades, and extensive external tubing for fuel, oil and air systems; they use temperature resistant alloys, and are made with tight specifications requiring precision manufacture. All this often makes the construction of a simple gas turbine more complicated than a piston engine.
Moreover, to reach optimum performance in modern gas turbine power plants the gas needs to be prepared to exact fuel specifications. Fuel gas conditioning systems treat the natural gas to reach the exact fuel specification prior to entering the turbine in terms of pressure, temperature, gas composition, and the related Wobbe index.
The primary advantage of a gas turbine engine is its power to weight ratio.
Since significant useful work can be generated by a relatively lightweight engine, gas turbines are perfectly suited for aircraft propulsion.
Thrust bearings and journal bearings are a critical part of a design. They are hydrodynamic oil bearings or oil-cooled rolling-element bearings. Foil bearings are used in some small machines such as micro turbines and also have strong potential for use in small gas turbines/auxiliary power units
Creep
A major challenge facing turbine design, especially turbine blades, is reducing the creep that is induced by the high temperatures and stresses that are experienced during operation. Higher operating temperatures are continuously sought in order to increase efficiency, but come at the cost of higher creep rates. Several methods have therefore been employed in an attempt to achieve optimal performance while limiting creep, with the most successful ones being high performance coatings and single crystal superalloys. These technologies work by limiting deformation that occurs by mechanisms that can be broadly classified as dislocation glide, dislocation climb and diffusional flow.
Protective coatings provide thermal insulation of the blade and offer oxidation and corrosion resistance. Thermal barrier coatings (TBCs) are often stabilized zirconium dioxide-based ceramics and oxidation/corrosion resistant coatings (bond coats) typically consist of aluminides or MCrAlY (where M is typically Fe and/or Cr) alloys. Using TBCs limits the temperature exposure of the superalloy substrate, thereby decreasing the diffusivity of the active species (typically vacancies) within the alloy and reducing dislocation and vacancy creep. It has been found that a coating of 1–200 μm can decrease blade temperatures by up to .
Bond coats are directly applied onto the surface of the substrate using pack carburization and serve the dual purpose of providing improved adherence for the TBC and oxidation resistance for the substrate. The Al from the bond coats forms Al2O3 on the TBC-bond coat interface which provides the oxidation resistance, but also results in the formation of an undesirable interdiffusion (ID) zone between itself and the substrate. The oxidation resistance outweighs the drawbacks associated with the ID zone as it increases the lifetime of the blade and limits the efficiency losses caused by a buildup on the outside of the blades.
Nickel-based superalloys boast improved strength and creep resistance due to their composition and resultant microstructure. The gamma (γ) FCC nickel is alloyed with aluminum and titanium in order to precipitate a uniform dispersion of the coherent gamma-prime (γ') phases. The finely dispersed γ' precipitates impede dislocation motion and introduce a threshold stress, increasing the stress required for the onset of creep. Furthermore, γ' is an ordered L12 phase that makes it harder for dislocations to shear past it. Further Refractory elements such as rhenium and ruthenium can be added in solid solution to improve creep strength. The addition of these elements reduces the diffusion of the gamma prime phase, thus preserving the fatigue resistance, strength, and creep resistance. The development of single crystal superalloys has led to significant improvements in creep resistance as well. Due to the lack of grain boundaries, single crystals eliminate Coble creep and consequently deform by fewer modes – decreasing the creep rate. Although single crystals have lower creep at high temperatures, they have significantly lower yield stresses at room temperature where strength is determined by the Hall-Petch relationship. Care needs to be taken in order to optimize the design parameters to limit high temperature creep while not decreasing low temperature yield strength.
Types
Jet engines
Airbreathing jet engines are gas turbines optimized to produce thrust from the exhaust gases, or from ducted fans connected to the gas turbines. Jet engines that produce thrust from the direct impulse of exhaust gases are often called turbojets. While still in service with many militaries and civilian operators, turbojets have mostly been phased out in favor of the turbofan engine due to the turbojet's low fuel efficiency, and high noise. Those that generate thrust with the addition of a ducted fan are called turbofans or (rarely) fan-jets. These engines produce nearly 80% of their thrust by the ducted fan, which can be seen from the front of the engine. They come in two types, low-bypass turbofan and high bypass, the difference being the amount of air moved by the fan, called "bypass air". These engines offer the benefit of more thrust without extra fuel consumption.
Gas turbines are also used in many liquid-fuel rockets, where gas turbines are used to power a turbopump to permit the use of lightweight, low-pressure tanks, reducing the empty weight of the rocket.
Turboprop engines
A turboprop engine is a turbine engine that drives an aircraft propeller using a reduction gear to translate high turbine section operating speed (often in the 10s of thousands) into low thousands necessary for efficient propeller operation. The benefit of using the turboprop engine is to take advantage of the turbine engines high power-to-weight ratio to drive a propeller, thus allowing a more powerful, but also smaller engine to be used. Turboprop engines are used on a wide range of business aircraft such as the Pilatus PC-12, commuter aircraft such as the Beechcraft 1900, and small cargo aircraft such as the Cessna 208 Caravan or De Havilland Canada Dash 8, and large aircraft (typically military) such as the Airbus A400M transport, Lockheed AC-130 and the 60-year-old Tupolev Tu-95 strategic bomber. While military turboprop engines can vary, in the civilian market there are two primary engines to be found: the Pratt & Whitney Canada PT6, a free-turbine turboshaft engine, and the Honeywell TPE331, a fixed turbine engine (formerly designated as the Garrett AiResearch 331).
Aeroderivative gas turbines
Aeroderivative gas turbines are generally based on existing aircraft gas turbine engines and are smaller and lighter than industrial gas turbines.
Aeroderivatives are used in electrical power generation due to their ability to be shut down and handle load changes more quickly than industrial machines. They are also used in the marine industry to reduce weight. Common types include the General Electric LM2500, General Electric LM6000, and aeroderivative versions of the Pratt & Whitney PW4000, Pratt & Whitney FT4 and Rolls-Royce RB211.
Amateur gas turbines
Increasing numbers of gas turbines are being used or even constructed by amateurs.
In its most straightforward form, these are commercial turbines acquired through military surplus or scrapyard sales, then operated for display as part of the hobby of engine collecting. In its most extreme form, amateurs have even rebuilt engines beyond professional repair and then used them to compete for the land speed record.
The simplest form of self-constructed gas turbine employs an automotive turbocharger as the core component. A combustion chamber is fabricated and plumbed between the compressor and turbine sections.
More sophisticated turbojets are also built, where their thrust and light weight are sufficient to power large model aircraft. The Schreckling design constructs the entire engine from raw materials, including the fabrication of a centrifugal compressor wheel from plywood, epoxy and wrapped carbon fibre strands.
Several small companies now manufacture small turbines and parts for the amateur. Most turbojet-powered model aircraft are now using these commercial and semi-commercial microturbines, rather than a Schreckling-like home-build.
Auxiliary power units
Small gas turbines are used as auxiliary power units (APUs) to supply auxiliary power to larger, mobile, machines such as an aircraft, and are a turboshaft design. They supply:
compressed air for air cycle machine style air conditioning and ventilation,
compressed air start-up power for larger jet engines,
mechanical (shaft) power to a gearbox to drive shafted accessories, and
electrical, hydraulic and other power-transmission sources to consuming devices remote from the APU.
Industrial gas turbines for power generation
Industrial gas turbines differ from aeronautical designs in that the frames, bearings, and blading are of heavier construction. They are also much more closely integrated with the devices they power—often an electric generator—and the secondary-energy equipment that is used to recover residual energy (largely heat).
They range in size from portable mobile plants to large, complex systems weighing more than a hundred tonnes housed in purpose-built buildings. When the gas turbine is used solely for shaft power, its thermal efficiency is about 30%. However, it may be cheaper to buy electricity than to generate it. Therefore, many engines are used in CHP (Combined Heat and Power) configurations that can be small enough to be integrated into portable container configurations.
Gas turbines can be particularly efficient when waste heat from the turbine is recovered by a heat recovery steam generator (HRSG) to power a conventional steam turbine in a combined cycle configuration. The 605 MW General Electric 9HA achieved a 62.22% efficiency rate with temperatures as high as .
For 2018, GE offers its 826 MW HA at over 64% efficiency in combined cycle due to advances in additive manufacturing and combustion breakthroughs, up from 63.7% in 2017 orders and on track to achieve 65% by the early 2020s.
In March 2018, GE Power achieved a 63.08% gross efficiency for its 7HA turbine.
Aeroderivative gas turbines can also be used in combined cycles, leading to a higher efficiency, but it will not be as high as a specifically designed industrial gas turbine. They can also be run in a cogeneration configuration: the exhaust is used for space or water heating, or drives an absorption chiller for cooling the inlet air and increase the power output, technology known as turbine inlet air cooling.
Another significant advantage is their ability to be turned on and off within minutes, supplying power during peak, or unscheduled, demand. Since single cycle (gas turbine only) power plants are less efficient than combined cycle plants, they are usually used as peaking power plants, which operate anywhere from several hours per day to a few dozen hours per year—depending on the electricity demand and the generating capacity of the region. In areas with a shortage of base-load and load following power plant capacity or with low fuel costs, a gas turbine powerplant may regularly operate most hours of the day. A large single-cycle gas turbine typically produces 100 to 400 megawatts of electric power and has 35–40% thermodynamic efficiency.
Industrial gas turbines for mechanical drive
Industrial gas turbines that are used solely for mechanical drive or used in collaboration with a recovery steam generator differ from power generating sets in that they are often smaller and feature a dual shaft design as opposed to a single shaft. The power range varies from 1 megawatt up to 50 megawatts. These engines are connected directly or via a gearbox to either a pump or compressor assembly. The majority of installations are used within the oil and gas industries. Mechanical drive applications increase efficiency by around 2%.
Oil and gas platforms require these engines to drive compressors to inject gas into the wells to force oil up via another bore, or to compress the gas for transportation. They are also often used to provide power for the platform. These platforms do not need to use the engine in collaboration with a CHP system due to getting the gas at an extremely reduced cost (often free from burn off gas). The same companies use pump sets to drive the fluids to land and across pipelines in various intervals.
Compressed air energy storage
One modern development seeks to improve efficiency in another way, by separating the compressor and the turbine with a compressed air store. In a conventional turbine, up to half the generated power is used driving the compressor. In a compressed air energy storage configuration, power is used to drive the compressor, and the compressed air is released to operate the turbine when required.
Turboshaft engines
Turboshaft engines are used to drive compressors in gas pumping stations and natural gas liquefaction plants. They are also used in aviation to power all but the smallest modern helicopters, and function as an auxiliary power unit in large commercial aircraft. A primary shaft carries the compressor and its turbine which, together with a combustor, is called a Gas Generator. A separately spinning power-turbine is usually used to drive the rotor on helicopters. Allowing the gas generator and power turbine/rotor to spin at their own speeds allows more flexibility in their design.
Radial gas turbines
Scale jet engines
Also known as miniature gas turbines or micro-jets.
With this in mind the pioneer of modern Micro-Jets, Kurt Schreckling, produced one of the world's first Micro-Turbines, the FD3/67. This engine can produce up to 22 newtons of thrust, and can be built by most mechanically minded people with basic engineering tools, such as a metal lathe.
Microturbines
Evolved from piston engine turbochargers, aircraft APUs or small jet engines, microturbines are 25 to 500 kilowatt turbines the size of a refrigerator.
Microturbines have around 15% efficiencies without a recuperator, 20 to 30% with one and they can reach 85% combined thermal-electrical efficiency in cogeneration.
External combustion
Most gas turbines are internal combustion engines but it is also possible to manufacture an external combustion gas turbine which is, effectively, a turbine version of a hot air engine.
Those systems are usually indicated as EFGT (Externally Fired Gas Turbine) or IFGT (Indirectly Fired Gas Turbine).
External combustion has been used for the purpose of using pulverized coal or finely ground biomass (such as sawdust) as a fuel. In the indirect system, a heat exchanger is used and only clean air with no combustion products travels through the power turbine. The thermal efficiency is lower in the indirect type of external combustion; however, the turbine blades are not subjected to combustion products and much lower quality (and therefore cheaper) fuels are able to be used.
When external combustion is used, it is possible to use exhaust air from the turbine as the primary combustion air. This effectively reduces global heat losses, although heat losses associated with the combustion exhaust remain inevitable.
Closed-cycle gas turbines based on helium or supercritical carbon dioxide also hold promise for use with future high temperature solar and nuclear power generation.
In surface vehicles
Gas turbines are often used on ships, locomotives, helicopters, tanks, and to a lesser extent, on cars, buses, and motorcycles.
A key advantage of jets and turboprops for airplane propulsion – their superior performance at high altitude compared to piston engines, particularly naturally aspirated ones – is irrelevant in most automobile applications. Their power-to-weight advantage, though less critical than for aircraft, is still important.
Gas turbines offer a high-powered engine in a very small and light package. However, they are not as responsive and efficient as small piston engines over the wide range of RPMs and powers needed in vehicle applications. In series hybrid vehicles, as the driving electric motors are mechanically detached from the electricity generating engine, the responsiveness, poor performance at low speed and low efficiency at low output problems are much less important. The turbine can be run at optimum speed for its power output, and batteries and ultracapacitors can supply power as needed, with the engine cycled on and off to run it only at high efficiency. The emergence of the continuously variable transmission may also alleviate the responsiveness problem.
Turbines have historically been more expensive to produce than piston engines, though this is partly because piston engines have been mass-produced in huge quantities for decades, while small gas turbine engines are rarities; however, turbines are mass-produced in the closely related form of the turbocharger.
The turbocharger is basically a compact and simple free shaft radial gas turbine which is driven by the piston engine's exhaust gas. The centripetal turbine wheel drives a centrifugal compressor wheel through a common rotating shaft. This wheel supercharges the engine air intake to a degree that can be controlled by means of a wastegate or by dynamically modifying the turbine housing's geometry (as in a variable geometry turbocharger).
It mainly serves as a power recovery device which converts a great deal of otherwise wasted thermal and kinetic energy into engine boost.
Turbo-compound engines (actually employed on some semi-trailer trucks) are fitted with blow down turbines which are similar in design and appearance to a turbocharger except for the turbine shaft being mechanically or hydraulically connected to the engine's crankshaft instead of to a centrifugal compressor, thus providing additional power instead of boost. While the turbocharger is a pressure turbine, a power recovery turbine is a velocity one.
Passenger road vehicles (cars, bikes, and buses)
A number of experiments have been conducted with gas turbine powered automobiles, the largest by Chrysler. More recently, there has been some interest in the use of turbine engines for hybrid electric cars. For instance, a consortium led by micro gas turbine company Bladon Jets has secured investment from the Technology Strategy Board to develop an Ultra Lightweight Range Extender (ULRE) for next-generation electric vehicles. The objective of the consortium, which includes luxury car maker Jaguar Land Rover and leading electrical machine company SR Drives, is to produce the world's first commercially viable – and environmentally friendly – gas turbine generator designed specifically for automotive applications.
The common turbocharger for gasoline or diesel engines is also a turbine derivative.
Concept cars
The first serious investigation of using a gas turbine in cars was in 1946 when two engineers, Robert Kafka and Robert Engerstein of Carney Associates, a New York engineering firm, came up with the concept where a unique compact turbine engine design would provide power for a rear wheel drive car. After an article appeared in Popular Science, there was no further work, beyond the paper stage.
Early concepts (1950s/60s)
In 1950, designer F.R. Bell and Chief Engineer Maurice Wilks from British car manufacturers Rover unveiled the first car powered with a gas turbine engine. The two-seater JET1 had the engine positioned behind the seats, air intake grilles on either side of the car, and exhaust outlets on the top of the tail. During tests, the car reached top speeds of , at a turbine speed of 50,000 rpm. After being shown in the United Kingdom and the United States in 1950, JET1 was further developed, and was subjected to speed trials on the Jabbeke highway in Belgium in June 1952, where it exceeded . The car ran on petrol, paraffin (kerosene) or diesel oil, but fuel consumption problems proved insurmountable for a production car. JET1 is on display at the London Science Museum.
A French turbine-powered car, the SOCEMA-Grégoire, was displayed at the October 1952 Paris Auto Show. It was designed by the French engineer Jean-Albert Grégoire.
The first turbine-powered car built in the US was the GM Firebird I which began evaluations in 1953. While photos of the Firebird I may suggest that the jet turbine's thrust propelled the car like an aircraft, the turbine actually drove the rear wheels. The Firebird I was never meant as a commercial passenger car and was built solely for testing & evaluation as well as public relation purposes. Additional Firebird concept cars, each powered by gas turbines, were developed for the 1953, 1956 and 1959 Motorama auto shows. The GM Research gas turbine engine also was fitted to a series of transit buses, starting with the Turbo-Cruiser I of 1953.
Starting in 1954 with a modified Plymouth, the American car manufacturer Chrysler demonstrated several prototype gas turbine-powered cars from the early 1950s through the early 1980s. Chrysler built fifty Chrysler Turbine Cars in 1963 and conducted the only consumer trial of gas turbine-powered cars. Each of their turbines employed a unique rotating recuperator, referred to as a regenerator that increased efficiency.
In 1954, Fiat unveiled a concept car with a turbine engine, called Fiat Turbina. This vehicle, looking like an aircraft with wheels, used a unique combination of both jet thrust and the engine driving the wheels. Speeds of were claimed.
In the 1960s, Ford and GM also were developing gas turbine semi-trucks. Ford displayed the Big Red at the 1964 World's Fair. With the trailer, it was long, high, and painted crimson red. It contained the Ford-developed gas turbine engine, with output power and torque of and . The cab boasted a highway map of the continental U.S., a mini-kitchen, bathroom, and a TV for the co-driver. The fate of the truck was unknown for several decades, but it was rediscovered in early 2021 in private hands, having been restored to running order. The Chevrolet division of GM built the Turbo Titan series of concept trucks with turbine motors as analogs of the Firebird concepts, including Turbo Titan I (, shares GT-304 engine with Firebird II), Turbo Titan II (, shares GT-305 engine with Firebird III), and Turbo Titan III (1965, GT-309 engine); in addition, the GM Bison gas turbine truck was shown at the 1964 World's Fair.
Emissions and fuel economy (1970s/80s)
As a result of the U.S. Clean Air Act Amendments of 1970, research was funded into developing automotive gas turbine technology. Design concepts and vehicles were conducted by Chrysler, General Motors, Ford (in collaboration with AiResearch), and American Motors (in conjunction with Williams Research). Long-term tests were conducted to evaluate comparable cost efficiency. Several AMC Hornets were powered by a small Williams regenerative gas turbine weighing and producing at 4450 rpm.
In 1982, General Motors used an Oldsmobile Delta 88 powered by a gas turbine using pulverised coal dust. This was considered for the United States and the western world to reduce dependence on middle east oil at the time
Toyota demonstrated several gas turbine powered concept cars, such as the Century gas turbine hybrid in 1975, the Sports 800 Gas Turbine Hybrid in 1979 and the GTV in 1985. No production vehicles were made. The GT24 engine was exhibited in 1977 without a vehicle.
Later development
In the early 1990s, Volvo introduced the Volvo ECC which was a gas turbine powered hybrid electric vehicle.
In 1993, General Motors developed a gas turbine powered EV1 series hybrid—as a prototype of the General Motors EV1. A Williams International 40 kW turbine drove an alternator which powered the battery–electric powertrain. The turbine design included a recuperator. In 2006, GM went into the EcoJet concept car project with Jay Leno.
At the 2010 Paris Motor Show Jaguar demonstrated its Jaguar C-X75 concept car. This electrically powered supercar has a top speed of and can go from in 3.4 seconds. It uses lithium-ion batteries to power four electric motors which combine to produce 780 bhp. It will travel on a single charge of the batteries, and uses a pair of Bladon Micro Gas Turbines to re-charge the batteries extending the range to .
Racing cars
The first race car (in concept only) fitted with a turbine was in 1955 by a US Air Force group as a hobby project with a turbine loaned them by Boeing and a race car owned by Firestone Tire & Rubber company. The first race car fitted with a turbine for the goal of actual racing was by Rover and the BRM Formula One team joined forces to produce the Rover-BRM, a gas turbine powered coupe, which entered the 1963 24 Hours of Le Mans, driven by Graham Hill and Richie Ginther. It averaged and had a top speed of . American Ray Heppenstall joined Howmet Corporation and McKee Engineering together to develop their own gas turbine sports car in 1968, the Howmet TX, which ran several American and European events, including two wins, and also participated in the 1968 24 Hours of Le Mans. The cars used Continental gas turbines, which eventually set six FIA land speed records for turbine-powered cars.
For open wheel racing, 1967's revolutionary STP-Paxton Turbocar fielded by racing and entrepreneurial legend Andy Granatelli and driven by Parnelli Jones nearly won the Indianapolis 500; the Pratt & Whitney ST6B-62 powered turbine car was almost a lap ahead of the second place car when a gearbox bearing failed just three laps from the finish line. The next year the STP Lotus 56 turbine car won the Indianapolis 500 pole position even though new rules restricted the air intake dramatically. In 1971 Team Lotus principal Colin Chapman introduced the Lotus 56B F1 car, powered by a Pratt & Whitney STN 6/76 gas turbine. Chapman had a reputation of building radical championship-winning cars, but had to abandon the project because there were too many problems with turbo lag.
Buses
General Motors fitted the GT-30x series of gas turbines (branded "Whirlfire") to several prototype buses in the 1950s and 1960s, including Turbo-Cruiser I (1953, GT-300); Turbo-Cruiser II (1964, GT-309); Turbo-Cruiser III (1968, GT-309); RTX (1968, GT-309); and RTS 3T (1972).
The arrival of the Capstone Turbine has led to several hybrid bus designs, starting with HEV-1 by AVS of Chattanooga, Tennessee in 1999, and closely followed by Ebus and ISE Research in California, and DesignLine Corporation in New Zealand (and later the United States). AVS turbine hybrids were plagued with reliability and quality control problems, resulting in liquidation of AVS in 2003. The most successful design by Designline is now operated in 5 cities in 6 countries, with over 30 buses in operation worldwide, and order for several hundred being delivered to Baltimore, and New York City.
Brescia Italy is using serial hybrid buses powered by microturbines on routes through the historical sections of the city.
Motorcycles
The MTT Turbine Superbike appeared in 2000 (hence the designation of Y2K Superbike by MTT) and is the first production motorcycle powered by a turbine engine – specifically, a Rolls-Royce Allison model 250 turboshaft engine, producing about 283 kW (380 bhp). Speed-tested to 365 km/h or 227 mph (according to some stories, the testing team ran out of road during the test), it holds the Guinness World Record for most powerful production motorcycle and most expensive production motorcycle, with a price tag of US$185,000.
Trains
Several locomotive classes have been powered by gas turbines, the most recent incarnation being Bombardier's JetTrain.
Tanks
The Third Reich Wehrmacht Heer's development division, the Heereswaffenamt (Army Ordnance Board), studied a number of gas turbine engine designs for use in tanks starting in mid-1944. The first gas turbine engine design intended for use in armored fighting vehicle propulsion, the BMW 003-based GT 101, was meant for installation in the Panther tank. Towards the end of the war, a Jagdtiger was fitted with one of the aforementioned gas turbines.
The second use of a gas turbine in an armored fighting vehicle was in 1954 when a unit, PU2979, specifically developed for tanks by C. A. Parsons and Company, was installed and trialed in a British Conqueror tank. The Stridsvagn 103 was developed in the 1950s and was the first mass-produced main battle tank to use a turbine engine, the Boeing T50. Since then, gas turbine engines have been used as auxiliary power units in some tanks and as main powerplants in Soviet/Russian T-80s and U.S. M1 Abrams tanks, among others. They are lighter and smaller than diesel engines at the same sustained power output but the models installed to date are less fuel efficient than the equivalent diesel, especially at idle, requiring more fuel to achieve the same combat range. Successive models of M1 have addressed this problem with battery packs or secondary generators to power the tank's systems while stationary, saving fuel by reducing the need to idle the main turbine. T-80s can mount three large external fuel drums to extend their range. Russia has stopped production of the T-80 in favor of the diesel-powered T-90 (based on the T-72), while Ukraine has developed the diesel-powered T-80UD and T-84 with nearly the power of the gas-turbine tank. The French Leclerc tank's diesel powerplant features the "Hyperbar" hybrid supercharging system, where the engine's turbocharger is completely replaced with a small gas turbine which also works as an assisted diesel exhaust turbocharger, enabling engine RPM-independent boost level control and a higher peak boost pressure to be reached (than with ordinary turbochargers). This system allows a smaller displacement and lighter engine to be used as the tank's power plant and effectively removes turbo lag. This special gas turbine/turbocharger can also work independently from the main engine as an ordinary APU.
A turbine is theoretically more reliable and easier to maintain than a piston engine since it has a simpler construction with fewer moving parts, but in practice, turbine parts experience a higher wear rate due to their higher working speeds. The turbine blades are highly sensitive to dust and fine sand so that in desert operations air filters have to be fitted and changed several times daily. An improperly fitted filter, or a bullet or shell fragment that punctures the filter, can damage the engine. Piston engines (especially if turbocharged) also need well-maintained filters, but they are more resilient if the filter does fail.
Like most modern diesel engines used in tanks, gas turbines are usually multi-fuel engines.
Marine applications
Naval
Gas turbines are used in many naval vessels, where they are valued for their high power-to-weight ratio and their ships' resulting acceleration and ability to get underway quickly.
The first gas-turbine-powered naval vessel was the Royal Navy's motor gunboat MGB 2009 (formerly MGB 509) converted in 1947. Metropolitan-Vickers fitted their F2/3 jet engine with a power turbine. The Steam Gun Boat Grey Goose was converted to Rolls-Royce gas turbines in 1952 and operated as such from 1953. The Bold class Fast Patrol Boats Bold Pioneer and Bold Pathfinder built in 1953 were the first ships created specifically for gas turbine propulsion.
The first large-scale, partially gas-turbine powered ships were the Royal Navy's Type 81 (Tribal class) frigates with combined steam and gas powerplants. The first, was commissioned in 1961.
The German Navy launched the first in 1961 with 2 Brown, Boveri & Cie gas turbines in the world's first combined diesel and gas propulsion system.
The Soviet Navy commissioned in 1962 the first of 25 with 4 gas turbines in combined gas and gas propulsion system. Those vessels used 4 M8E gas turbines, which generated . Those ships were the first large ships in the world to be powered solely by gas turbines.
The Danish Navy had 6 Søløven-class torpedo boats (the export version of the British Brave class fast patrol boat) in service from 1965 to 1990, which had 3 Bristol Proteus (later RR Proteus) Marine Gas Turbines rated at combined, plus two General Motors Diesel engines, rated at , for better fuel economy at slower speeds. And they also produced 10 Willemoes Class Torpedo / Guided Missile boats (in service from 1974 to 2000) which had 3 Rolls-Royce Marine Proteus Gas Turbines also rated at , same as the Søløven-class boats, and 2 General Motors Diesel Engines, rated at , also for improved fuel economy at slow speeds.
The Swedish Navy produced 6 Spica-class torpedo boats between 1966 and 1967 powered by 3 Bristol Siddeley Proteus 1282 turbines, each delivering . They were later joined by 12 upgraded Norrköping class ships, still with the same engines. With their aft torpedo tubes replaced by antishipping missiles they served as missile boats until the last was retired in 2005.
The Finnish Navy commissioned two corvettes, Turunmaa and Karjala, in 1968. They were equipped with one Rolls-Royce Olympus TM1 gas turbine and three Wärtsilä marine diesels for slower speeds. They were the fastest vessels in the Finnish Navy; they regularly achieved speeds of 35 knots, and 37.3 knots during sea trials. The Turunmaas were decommissioned in 2002. Karjala is today a museum ship in Turku, and Turunmaa serves as a floating machine shop and training ship for Satakunta Polytechnical College.
The next series of major naval vessels were the four Canadian helicopter carrying destroyers first commissioned in 1972. They used 2 ft-4 main propulsion engines, 2 ft-12 cruise engines and 3 Solar Saturn 750 kW generators.
The first U.S. gas-turbine powered ship was the U.S. Coast Guard's , a cutter commissioned in 1961 that was powered by two turbines utilizing controllable-pitch propellers. The larger High Endurance Cutters, was the first class of larger cutters to utilize gas turbines, the first of which () was commissioned in 1967. Since then, they have powered the U.S. Navy's s, and s, and guided missile cruisers. , a modified , is to be the Navy's first amphibious assault ship powered by gas turbines.
The marine gas turbine operates in a more corrosive atmosphere due to the presence of sea salt in air and fuel and use of cheaper fuels.
Civilian maritime
Up to the late 1940s, much of the progress on marine gas turbines all over the world took place in design offices and engine builder's workshops and development work was led by the British Royal Navy and other Navies. While interest in the gas turbine for marine purposes, both naval and mercantile, continued to increase, the lack of availability of the results of operating experience on early gas turbine projects limited the number of new ventures on seagoing commercial vessels being embarked upon.
In 1951, the diesel–electric oil tanker Auris, 12,290 deadweight tonnage (DWT) was used to obtain operating experience with a main propulsion gas turbine under service conditions at sea and so became the first ocean-going merchant ship to be powered by a gas turbine. Built by Hawthorn Leslie at Hebburn-on-Tyne, UK, in accordance with plans and specifications drawn up by the Anglo-Saxon Petroleum Company and launched on the UK's Princess Elizabeth's 21st birthday in 1947, the ship was designed with an engine room layout that would allow for the experimental use of heavy fuel in one of its high-speed engines, as well as the future substitution of one of its diesel engines by a gas turbine. The Auris operated commercially as a tanker for three-and-a-half years with a diesel–electric propulsion unit as originally commissioned, but in 1951 one of its four diesel engines – which were known as "Faith", "Hope", "Charity" and "Prudence" – was replaced by the world's first marine gas turbine engine, a open-cycle gas turbo-alternator built by British Thompson-Houston Company in Rugby. Following successful sea trials off the Northumbrian coast, the Auris set sail from Hebburn-on-Tyne in October 1951 bound for Port Arthur in the US and then Curaçao in the southern Caribbean returning to Avonmouth after 44 days at sea, successfully completing her historic trans-Atlantic crossing. During this time at sea the gas turbine burnt diesel fuel and operated without an involuntary stop or mechanical difficulty of any kind. She subsequently visited Swansea, Hull, Rotterdam, Oslo and Southampton covering a total of 13,211 nautical miles. The Auris then had all of its power plants replaced with a directly coupled gas turbine to become the first civilian ship to operate solely on gas turbine power.
Despite the success of this early experimental voyage the gas turbine did not replace the diesel engine as the propulsion plant for large merchant ships. At constant cruising speeds the diesel engine simply had no peer in the vital area of fuel economy. The gas turbine did have more success in Royal Navy ships and the other naval fleets of the world where sudden and rapid changes of speed are required by warships in action.
The United States Maritime Commission were looking for options to update WWII Liberty ships, and heavy-duty gas turbines were one of those selected. In 1956 the John Sergeant was lengthened and equipped with a General Electric HD gas turbine with exhaust-gas regeneration, reduction gearing and a variable-pitch propeller. It operated for 9,700 hours using residual fuel (Bunker C) for 7,000 hours. Fuel efficiency was on a par with steam propulsion at per hour, and power output was higher than expected at due to the ambient temperature of the North Sea route being lower than the design temperature of the gas turbine. This gave the ship a speed capability of 18 knots, up from 11 knots with the original power plant, and well in excess of the 15 knot targeted. The ship made its first transatlantic crossing with an average speed of 16.8 knots, in spite of some rough weather along the way. Suitable Bunker C fuel was only available at limited ports because the quality of the fuel was of a critical nature. The fuel oil also had to be treated on board to reduce contaminants and this was a labor-intensive process that was not suitable for automation at the time. Ultimately, the variable-pitch propeller, which was of a new and untested design, ended the trial, as three consecutive annual inspections revealed stress-cracking. This did not reflect poorly on the marine-propulsion gas-turbine concept though, and the trial was a success overall. The success of this trial opened the way for more development by GE on the use of HD gas turbines for marine use with heavy fuels. The John Sergeant was scrapped in 1972 at Portsmouth PA.
Boeing launched its first passenger-carrying waterjet-propelled hydrofoil Boeing 929, in April 1974. Those ships were powered by two Allison 501-KF gas turbines.
Between 1971 and 1981, Seatrain Lines operated a scheduled container service between ports on the eastern seaboard of the United States and ports in northwest Europe across the North Atlantic with four container ships of 26,000 tonnes DWT. Those ships were powered by twin Pratt & Whitney gas turbines of the FT 4 series. The four ships in the class were named Euroliner, Eurofreighter, Asialiner and Asiafreighter. Following the dramatic Organization of the Petroleum Exporting Countries (OPEC) price increases of the mid-1970s, operations were constrained by rising fuel costs. Some modification of the engine systems on those ships was undertaken to permit the burning of a lower grade of fuel (i.e., marine diesel). Reduction of fuel costs was successful using a different untested fuel in a marine gas turbine but maintenance costs increased with the fuel change. After 1981 the ships were sold and refitted with, what at the time, was more economical diesel-fueled engines but the increased engine size reduced cargo space.
The first passenger ferry to use a gas turbine was the GTS Finnjet, built in 1977 and powered by two Pratt & Whitney FT 4C-1 DLF turbines, generating and propelling the ship to a speed of 31 knots. However, the Finnjet also illustrated the shortcomings of gas turbine propulsion in commercial craft, as high fuel prices made operating her unprofitable. After four years of service, additional diesel engines were installed on the ship to reduce running costs during the off-season. The Finnjet was also the first ship with a combined diesel–electric and gas propulsion. Another example of commercial use of gas turbines in a passenger ship is Stena Line's HSS class fastcraft ferries. HSS 1500-class Stena Explorer, Stena Voyager and Stena Discovery vessels use combined gas and gas setups of twin GE LM2500 plus GE LM1600 power for a total of . The slightly smaller HSS 900-class Stena Carisma, uses twin ABB–STAL GT35 turbines rated at gross. The Stena Discovery was withdrawn from service in 2007, another victim of too high fuel costs.
In July 2000, the Millennium became the first cruise ship to be powered by both gas and steam turbines. The ship featured two General Electric LM2500 gas turbine generators whose exhaust heat was used to operate a steam turbine generator in a COGES (combined gas electric and steam) configuration. Propulsion was provided by two electrically driven Rolls-Royce Mermaid azimuth pods. The liner uses a combined diesel and gas configuration.
In marine racing applications the 2010 C5000 Mystic catamaran Miss GEICO uses two Lycoming T-55 turbines for its power system.
Advances in technology
Gas turbine technology has steadily advanced since its inception and continues to evolve. Development is actively producing both smaller gas turbines and more powerful and efficient engines. Aiding in these advances are computer-based design (specifically computational fluid dynamics and finite element analysis) and the development of advanced materials: Base materials with superior high-temperature strength (e.g., single-crystal superalloys that exhibit yield strength anomaly) or thermal barrier coatings that protect the structural material from ever-higher temperatures. These advances allow higher compression ratios and turbine inlet temperatures, more efficient combustion and better cooling of engine parts.
Computational fluid dynamics (CFD) has contributed to substantial improvements in the performance and efficiency of gas turbine engine components through enhanced understanding of the complex viscous flow and heat transfer phenomena involved. For this reason, CFD is one of the key computational tools used in design and development of gas turbine engines.
The simple-cycle efficiencies of early gas turbines were practically doubled by incorporating inter-cooling, regeneration (or recuperation), and reheating. These improvements, of course, come at the expense of increased initial and operation costs, and they cannot be justified unless the decrease in fuel costs offsets the increase in other costs. The relatively low fuel prices, the general desire in the industry to minimize installation costs, and the tremendous increase in the simple-cycle efficiency to about 40 percent left little desire for opting for these modifications.
On the emissions side, the challenge is to increase turbine inlet temperatures while at the same time reducing peak flame temperature in order to achieve lower NOx emissions and meet the latest emission regulations. In May 2011, Mitsubishi Heavy Industries achieved a turbine inlet temperature of on a 320 megawatt gas turbine, and 460 MW in gas turbine combined-cycle power generation applications in which gross thermal efficiency exceeds 60%.
Compliant foil bearings were commercially introduced to gas turbines in the 1990s. These can withstand over a hundred thousand start/stop cycles and have eliminated the need for an oil system. The application of microelectronics and power switching technology have enabled the development of commercially viable electricity generation by microturbines for distribution and vehicle propulsion.
In 2013, General Electric started the development of the GE9X with a compression ratio of 61:1.
Advantages and disadvantages
The following are advantages and disadvantages of gas-turbine engines:
Advantages include:
Very high power-to-weight ratio compared to reciprocating engines.
Smaller than most reciprocating engines of the same power rating.
Smooth rotation of the main shaft produces far less vibration than a reciprocating engine.
Fewer moving parts than reciprocating engines results in lower maintenance cost and higher reliability/availability over its service life.
Greater reliability, particularly in applications where sustained high power output is required.
Waste heat is dissipated almost entirely in the exhaust. This results in a high-temperature exhaust stream that is very usable for boiling water in a combined cycle, or for cogeneration.
Lower peak combustion pressures than reciprocating engines in general.
High shaft speeds in smaller "free turbine units", although larger gas turbines employed in power generation operate at synchronous speeds.
Low lubricating oil cost and consumption.
Can run on a wide variety of fuels.
Very low toxic emissions of CO and HC due to excess air, complete combustion and no "quench" of the flame on cold surfaces.
Disadvantages include:
Core engine costs can be high due to the use of exotic materials, especially in applications where high reliability is required (e.g. aircraft propulsion)
Less efficient than reciprocating engines at idle speed.
Longer startup than reciprocating engines.
Less responsive to changes in power demand compared with reciprocating engines.
Characteristic whine can be hard to suppress. The exhaust (particularly on turbojets) can also produce a distinctive roaring sound.
Major manufacturers
Siemens Energy
Ansaldo
Mitsubishi Heavy Industries
Rolls-Royce
GE Aviation
Silmash
ODK
Pratt & Whitney
P&W Canada
Solar Turbines
Alstom
Zorya-Mashproekt
MTU Aero Engines
MAN Turbo
IHI Corporation
Kawasaki Heavy Industries
HAL
BHEL
MAPNA
Techwin
Doosan Heavy
Shanghai Electric
Harbin Electric
AECC
Testing
British, German, other national and international test codes are used to standardize the procedures and definitions used to test gas turbines. Selection of the test code to be used is an agreement between the purchaser and the manufacturer, and has some significance to the design of the turbine and associated systems. In the United States, ASME has produced several performance test codes on gas turbines. This includes ASME PTC 22–2014. These ASME performance test codes have gained international recognition and acceptance for testing gas turbines. The single most important and differentiating characteristic of ASME performance test codes, including PTC 22, is that the test uncertainty of the measurement indicates the quality of the test and is not to be used as a commercial tolerance.
See also
List of aircraft engines
Centrifugal compressor
Gas turbine modular helium reactor
Pneumatic motor
Pulsejet
Steam turbine
Turbine engine failure
Wind turbine
References
Further reading
Stationary Combustion Gas Turbines including Oil & Over-Speed Control System description
"Aircraft Gas Turbine Technology" by Irwin E. Treager, McGraw-Hill, Glencoe Division, 1979, .
"Gas Turbine Theory" by H.I.H. Saravanamuttoo, G.F.C. Rogers and H. Cohen, Pearson Education, 2001, 5th ed., .
R. M. "Fred" Klaass and Christopher DellaCorte, "The Quest for Oil-Free Gas Turbine Engines," SAE Technical Papers, No. 2006-01-3055, available at sae.org
"Model Jet Engines" by Thomas Kamps Traplet Publications
Aircraft Engines and Gas Turbines, Second Edition by Jack L. Kerrebrock, The MIT Press, 1992, .
"Forensic Investigation of a Gas Turbine Event" by John Molloy, M&M Engineering
"Gas Turbine Performance, 2nd Edition" by Philip Walsh and Paul Fletcher, Wiley-Blackwell, 2004
External links
Technology Speed of Civil Jet Engines
MIT Gas Turbine Laboratory
MIT Microturbine research
California Distributed Energy Resource guide – Microturbine generators
Introduction to how a gas turbine works from "how stuff works.com"
Aircraft gas turbine simulator for interactive learning
An online handbook on stationary gas turbine technologies compiled by the US DOE.
Engines
Marine propulsion | Gas turbine | [
"Physics",
"Technology",
"Engineering"
] | 12,211 | [
"Machines",
"Engines",
"Gas turbines",
"Physical systems",
"Marine engineering",
"Marine propulsion"
] |
58,894 | https://en.wikipedia.org/wiki/Phosphorylation | In biochemistry, phosphorylation is the attachment of a phosphate group to a molecule or an ion. This process and its inverse, dephosphorylation, are common in biology. Protein phosphorylation often activates (or deactivates) many enzymes.
During respiration
Phosphorylation is essential to the processes of both anaerobic and aerobic respiration, which involve the production of adenosine triphosphate (ATP), the "high-energy" exchange medium in the cell. During aerobic respiration, ATP is synthesized in the mitochondrion by addition of a third phosphate group to adenosine diphosphate (ADP) in a process referred to as oxidative phosphorylation. ATP is also synthesized by substrate-level phosphorylation during glycolysis. ATP is synthesized at the expense of solar energy by photophosphorylation in the chloroplasts of plant cells.
Phosphorylation of glucose
Glucose metabolism
Phosphorylation of sugars is often the first stage in their catabolism. Phosphorylation allows cells to accumulate sugars because the phosphate group prevents the molecules from diffusing back across their transporter. Phosphorylation of glucose is a key reaction in sugar metabolism. The chemical equation for the conversion of D-glucose to D-glucose-6-phosphate in the first step of glycolysis is given by:
D-glucose + ATP → D-glucose 6-phosphate + ADP
ΔG° = −16.7 kJ/mol (° indicates measurement at standard condition)
Glycolysis
Glycolysis is an essential process of glucose degrading into two molecules of pyruvate, through various steps, with the help of different enzymes. It occurs in ten steps and proves that phosphorylation is a much required and necessary step to attain the end products. Phosphorylation initiates the reaction in step 1 of the preparatory step (first half of glycolysis), and initiates step 6 of payoff phase (second phase of glycolysis).
Glucose, by nature, is a small molecule with the ability to diffuse in and out of the cell. By phosphorylating glucose (adding a phosphoryl group in order to create a negatively charged phosphate group), glucose is converted to glucose-6-phosphate, which is trapped within the cell as the cell membrane is negatively charged. This reaction occurs due to the enzyme hexokinase, an enzyme that helps phosphorylate many six-membered ring structures. Phosphorylation takes place in step 3, where fructose-6-phosphate is converted to fructose 1,6-bisphosphate. This reaction is catalyzed by phosphofructokinase.
While phosphorylation is performed by ATPs during preparatory steps, phosphorylation during payoff phase is maintained by inorganic phosphate. Each molecule of glyceraldehyde 3-phosphate is phosphorylated to form 1,3-bisphosphoglycerate. This reaction is catalyzed by glyceraldehyde-3-phosphate dehydrogenase (GAPDH). The cascade effect of phosphorylation eventually causes instability and allows enzymes to open the carbon bonds in glucose.
Phosphorylation functions is an extremely vital component of glycolysis, as it helps in transport, control, and efficiency.
Glycogen synthesis
Glycogen is a long-term store of glucose produced by the cells of the liver. In the liver, the synthesis of glycogen is directly correlated with blood glucose concentration. High blood glucose concentration causes an increase in intracellular levels of glucose 6-phosphate in the liver, skeletal muscle, and fat (adipose) tissue. Glucose 6-phosphate has role in regulating glycogen synthase.
High blood glucose releases insulin, stimulating the translocation of specific glucose transporters to the cell membrane; glucose is phosphorylated to glucose 6-phosphate during transport across the membrane by ATP-D-glucose 6-phosphotransferase and non-specific hexokinase (ATP-D-hexose 6-phosphotransferase). Liver cells are freely permeable to glucose, and the initial rate of phosphorylation of glucose is the rate-limiting step in glucose metabolism by the liver.
The liver's crucial role in controlling blood sugar concentrations by breaking down glucose into carbon dioxide and glycogen is characterized by the negative Gibbs free energy (ΔG) value, which indicates that this is a point of regulation with. The hexokinase enzyme has a low Michaelis constant (K), indicating a high affinity for glucose, so this initial phosphorylation can proceed even when glucose levels at nanoscopic scale within the blood.
The phosphorylation of glucose can be enhanced by the binding of fructose 6-phosphate (F6P), and lessened by the binding fructose 1-phosphate (F1P). Fructose consumed in the diet is converted to F1P in the liver. This negates the action of F6P on glucokinase, which ultimately favors the forward reaction. The capacity of liver cells to phosphorylate fructose exceeds capacity to metabolize fructose-1-phosphate. Consuming excess fructose ultimately results in an imbalance in liver metabolism, which indirectly exhausts the liver cell's supply of ATP.
Allosteric activation by glucose-6-phosphate, which acts as an effector, stimulates glycogen synthase, and glucose-6-phosphate may inhibit the phosphorylation of glycogen synthase by cyclic AMP-stimulated protein kinase.
Other processes
Phosphorylation of glucose is imperative in processes within the body. For example, phosphorylating glucose is necessary for insulin-dependent mechanistic target of rapamycin pathway activity within the heart. This further suggests a link between intermediary metabolism and cardiac growth.
Protein phosphorylation
Protein phosphorylation is the most abundant post-translational modification in eukaryotes. Phosphorylation can occur on serine, threonine and tyrosine side chains (in other words, on their residues) through phosphoester bond formation, on histidine, lysine and arginine through phosphoramidate bonds, and on aspartic acid and glutamic acid through mixed anhydride linkages. Recent evidence confirms widespread histidine phosphorylation at both the 1 and 3 N-atoms of the imidazole ring. Recent work demonstrates widespread human protein phosphorylation on multiple non-canonical amino acids, including motifs containing phosphorylated histidine, aspartate, glutamate, cysteine, arginine and lysine in HeLa cell extracts. However, due to the chemical lability of these phosphorylated residues, and in marked contrast to Ser, Thr and Tyr phosphorylation, the analysis of phosphorylated histidine (and other non-canonical amino acids) using standard biochemical and mass spectrometric approaches is much more challenging and special procedures and separation techniques are required for their preservation alongside classical Ser, Thr and Tyr phosphorylation.
The prominent role of protein phosphorylation in biochemistry is illustrated by the huge body of studies published on the subject (as of March 2015, the MEDLINE database returns over 240,000 articles, mostly on protein phosphorylation).
See also
Moiety conservation
Phosida
Phosphoamino acid analysis
Phospho3D
References
External links
Functional analyses for site-specific phosphorylation of a target protein in cells (A Protocol)
Cell biology
Cell signaling
Phosphorus
Post-translational modification | Phosphorylation | [
"Chemistry",
"Biology"
] | 1,700 | [
"Post-translational modification",
"Gene expression",
"Cell biology",
"Biochemical reactions"
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59,046 | https://en.wikipedia.org/wiki/Silly%20Putty | Silly Putty is a toy containing silicone polymers that have unusual physical properties. It can flow like a liquid, bounce and can be stretched or broken depending on the amount of physical stress to which it is subjected. It contains viscoelastic liquid silicones, a type of non-Newtonian fluid, which makes it act as a viscous liquid over a long period of time but as an elastic solid over a short time period. It was originally created during research into a potential rubber substitute for use by the United States in World War II.
The name Silly Putty is a trademark of Crayola LLC. Other names are used to market similar substances from other manufacturers.
Description
As a bouncing putty, Silly Putty is noted for its unusual characteristics. It bounces when dropped from a height, but breaks when struck or stretched sharply; it can also float in a liquid and will form a puddle given enough time. Silly Putty and most other retail putty products have viscoelastic agents added to reduce the flow and enable the putty to hold its shape.
The original coral-colored Silly Putty is composed of 65% dimethylsiloxane (hydroxy-terminated polymers with boric acid), 17% silica (crystalline quartz), 9% Thixatrol ST (castor oil derivative), 4% polydimethylsiloxane, 1% decamethyl cyclopentasiloxane, 1% glycerine, and 1% titanium dioxide.
Silly Putty's unusual flow characteristics are due to the ingredient polydimethylsiloxane (PDMS), a viscoelastic substance. Viscoelasticity is a type of non-Newtonian flow, characterizing a material that acts as a viscous liquid over a long time period but as an elastic solid over a short time period. Because its apparent viscosity increases directly with respect to the amount of force applied, Silly Putty can be characterized as a dilatant fluid.
Silly Putty is also a fairly good adhesive. When newspaper ink was petroleum based, Silly Putty could be used to transfer newspaper images to other surfaces, providing amusement by distorting the transferred image afterwards. Newer papers with soy-based inks are more resistant to this process.
Generally, Silly Putty is difficult to remove from textured items such as dirt and clothing. Hand sanitizers containing alcohol are often helpful. Silly Putty will dissolve when in contact with an alcohol; after the alcohol evaporates, the material will not exhibit its original properties.
If Silly Putty is submerged in warm or hot water, it will become softer and thus "melt" much faster. It also becomes harder to remove small amounts of it from surfaces. After a long period of time, it will return to its original viscosity.
Silly Putty is sold as a piece of clay inside an egg-shaped plastic container. The Silly Putty brand is owned by Crayola LLC (formerly the Binney & Smith company). , twenty thousand eggs of Silly Putty are sold daily. Since 1950, more than 300 million eggs of Silly Putty (approximately ) have been sold. It is available in various colors, including glow-in-the-dark and metallic. Other brands offer similar materials, sometimes in larger-sized containers, and in a similarly wide variety of colors or with different properties, such as magnetism and iridescence.
History
During World War II, Japan invaded rubber-producing countries as it expanded its sphere of influence in the Pacific Rim. Rubber was vital for the production of rafts, tires, vehicle and aircraft parts, gas masks, and boots. In the US, all rubber products were rationed; citizens were encouraged to make their rubber products last until the end of the war and to donate spare tires, boots, and coats. Meanwhile, the government funded research into synthetic rubber compounds to attempt to solve this shortage.
Credit for the invention of Silly Putty is disputed and has been attributed variously to Earl Warrick of the then newly formed Dow Corning; Harvey Chin; and James Wright, a Scottish-born inventor working for General Electric in New Haven, Connecticut. Throughout his life, Warrick insisted that he and his colleague, Rob Roy McGregor, received the patent for Silly Putty before Wright did; but Crayola's history of Silly Putty states that Wright first invented it in 1943. Both researchers independently discovered that reacting boric acid with silicone oil would produce a gooey, bouncy material with several unique properties. The non-toxic putty would bounce when dropped, could stretch farther than regular rubber, would not go moldy, and had a very high melting temperature. However, the substance did not have all the properties needed to replace rubber.
In 1949, toy store owner Ruth Fallgatter came across the putty. She contacted marketing consultant Peter C. L. Hodgson (1912–1976). The two decided to market the bouncing putty by selling it in a clear case. Although it sold well, Fallgatter did not pursue it further. However, Hodgson saw its potential.
Already US$12,000 in debt, Hodgson borrowed $147 to buy a batch of the putty to pack portions into plastic eggs for $1, calling it Silly Putty. Initial sales were poor, but after a New Yorker article mentioned it, Hodgson sold over 250,000 eggs of silly putty in three days. However, Hodgson was almost put out of business in 1951 by the Korean War. Silicone, the main ingredient in silly putty, was put on ration, harming his business. A year later, the restriction on silicone was lifted and the production of Silly Putty resumed. Initially, it was primarily targeted towards adults. However, by 1955, the majority of its customers were aged six to twelve. In 1957, Hodgson produced the first televised commercial for Silly Putty, which aired during the Howdy Doody Show.
In 1961, Silly Putty went worldwide, becoming a hit in the Soviet Union and Europe. In 1968, it was taken into lunar orbit by the Apollo 8 astronauts.
Peter Hodgson died in 1976. A year later, Binney & Smith, the makers of Crayola products, acquired the rights to Silly Putty. , annual Silly Putty sales exceeded six million eggs.
Silly Putty was inducted into the National Toy Hall of Fame on May 28, 2001.
Other uses
In addition to its success as a toy, other uses for the putty have been found. In the home, it can be used to remove substances such as dirt, lint, pet hair, or ink from various surfaces. The material's unique properties have found niche use in medical and scientific applications. Occupational therapists use it for rehabilitative therapy of hand injuries. A number of other brands (such as Power Putty and TheraPutty) alter the material's properties, offering different levels of resistance. The material is also used as a tool to help reduce stress, and exists in various viscosities based on the user's preference.
Because of its adhesive characteristics, it was used by Apollo astronauts to secure their tools in zero gravity. Scale model building hobbyists use the putty as a masking medium when spray-painting model assemblies. The Steward Observatory uses a Silly-Putty backed lap to polish astronomical telescope mirrors.
Researchers from Trinity College Dublin School of Physics (Centre for Research on Adaptive Nanostructures and Nanodevices (CRANN) and Advanced Materials and Bioengineering Research (AMBER) Research Centers) have discovered nano composite mixtures of graphene and Silly Putty behave as sensitive pressure sensors, claiming the ability to measure the footsteps of a spider crawling on it.
See also
Blu Tack
Flubber (material)
Slime (toy)
References
External links
Clay toys
1940s toys
1950s toys
American inventions
Brand name materials
Crayola
Products introduced in 1949
Companies based in Northampton County, Pennsylvania
Dow Chemical Company
Easton, Pennsylvania
Non-Newtonian fluids
Polymers
Soft matter
Articles containing video clips
Sensory toys | Silly Putty | [
"Physics",
"Chemistry",
"Materials_science"
] | 1,646 | [
"Polymers",
"Soft matter",
"Condensed matter physics",
"Polymer chemistry"
] |
59,052 | https://en.wikipedia.org/wiki/Ensemble%20%28mathematical%20physics%29 | In physics, specifically statistical mechanics, an ensemble (also statistical ensemble) is an idealization consisting of a large number of virtual copies (sometimes infinitely many) of a system, considered all at once, each of which represents a possible state that the real system might be in. In other words, a statistical ensemble is a set of systems of particles used in statistical mechanics to describe a single
system. The concept of an ensemble was introduced by J. Willard Gibbs in 1902.
A thermodynamic ensemble is a specific variety of statistical ensemble that, among other properties, is in statistical equilibrium (defined below), and is used to derive the properties of thermodynamic systems from the laws of classical or quantum mechanics.
Physical considerations
The ensemble formalises the notion that an experimenter repeating an experiment again and again under the same macroscopic conditions, but unable to control the microscopic details, may expect to observe a range of different outcomes.
The notional size of ensembles in thermodynamics, statistical mechanics and quantum statistical mechanics can be very large, including every possible microscopic state the system could be in, consistent with its observed macroscopic properties. For many important physical cases, it is possible to calculate averages directly over the whole of the thermodynamic ensemble, to obtain explicit formulas for many of the thermodynamic quantities of interest, often in terms of the appropriate partition function.
The concept of an equilibrium or stationary ensemble is crucial to many applications of statistical ensembles. Although a mechanical system certainly evolves over time, the ensemble does not necessarily have to evolve. In fact, the ensemble will not evolve if it contains all past and future phases of the system. Such a statistical ensemble, one that does not change over time, is called stationary and can be said to be in statistical equilibrium.
Terminology
The word "ensemble" is also used for a smaller set of possibilities sampled from the full set of possible states. For example, a collection of walkers in a Markov chain Monte Carlo iteration is called an ensemble in some of the literature.
The term "ensemble" is often used in physics and the physics-influenced literature. In probability theory, the term probability space is more prevalent.
Main types
The study of thermodynamics is concerned with systems that appear to human perception to be "static" (despite the motion of their internal parts), and which can be described simply by a set of macroscopically observable variables. These systems can be described by statistical ensembles that depend on a few observable parameters, and which are in statistical equilibrium. Gibbs noted that different macroscopic constraints lead to different types of ensembles, with particular statistical characteristics.
"We may imagine a great number of systems of the same nature, but differing in the configurations and velocities which they have at a given instant, and differing in not merely infinitesimally, but it may be so as to embrace every conceivable combination of configuration and velocities..." J. W. Gibbs (1903)
Three important thermodynamic ensembles were defined by Gibbs:
Microcanonical ensemble (or NVE ensemble) —a statistical ensemble where the total energy of the system and the number of particles in the system are each fixed to particular values; each of the members of the ensemble are required to have the same total energy and particle number. The system must remain totally isolated (unable to exchange energy or particles with its environment) in order to stay in statistical equilibrium.
Canonical ensemble (or NVT ensemble)—a statistical ensemble where the energy is not known exactly but the number of particles is fixed. In place of the energy, the temperature is specified. The canonical ensemble is appropriate for describing a closed system which is in, or has been in, weak thermal contact with a heat bath. In order to be in statistical equilibrium, the system must remain totally closed (unable to exchange particles with its environment) and may come into weak thermal contact with other systems that are described by ensembles with the same temperature.
Grand canonical ensemble (or μVT ensemble)—a statistical ensemble where neither the energy nor particle number are fixed. In their place, the temperature and chemical potential are specified. The grand canonical ensemble is appropriate for describing an open system: one which is in, or has been in, weak contact with a reservoir (thermal contact, chemical contact, radiative contact, electrical contact, etc.). The ensemble remains in statistical equilibrium if the system comes into weak contact with other systems that are described by ensembles with the same temperature and chemical potential.
The calculations that can be made using each of these ensembles are explored further in their respective articles.
Other thermodynamic ensembles can be also defined, corresponding to different physical requirements, for which analogous formulae can often similarly be derived.
For example, in the reaction ensemble, particle number fluctuations are only allowed to occur according to the stoichiometry of the chemical reactions which are present in the system.
Equivalence
In thermodynamic limit all ensembles should produce identical observables due to Legendre transforms, deviations to this rule occurs under conditions that state-variables are non-convex, such as small molecular measurements.
Representations
The precise mathematical expression for a statistical ensemble has a distinct form depending on the type of mechanics under consideration (quantum or classical). In the classical case, the ensemble is a probability distribution over the microstates. In quantum mechanics, this notion, due to von Neumann, is a way of assigning a probability distribution over the results of each complete set of commuting observables.
In classical mechanics, the ensemble is instead written as a probability distribution in phase space; the microstates are the result of partitioning phase space into equal-sized units, although the size of these units can be chosen somewhat arbitrarily.
Requirements for representations
Putting aside for the moment the question of how statistical ensembles are generated operationally, we should be able to perform the following two operations on ensembles A, B of the same system:
Test whether A, B are statistically equivalent.
If p is a real number such that , then produce a new ensemble by probabilistic sampling from A with probability p and from B with probability .
Under certain conditions, therefore, equivalence classes of statistical ensembles have the structure of a convex set.
Quantum mechanical
A statistical ensemble in quantum mechanics (also known as a mixed state) is most often represented by a density matrix, denoted by . The density matrix provides a fully general tool that can incorporate both quantum uncertainties (present even if the state of the system were completely known) and classical uncertainties (due to a lack of knowledge) in a unified manner. Any physical observable in quantum mechanics can be written as an operator, . The expectation value of this operator on the statistical ensemble is given by the following trace:
This can be used to evaluate averages (operator ), variances (using operator ), covariances (using operator ), etc. The density matrix must always have a trace of 1: (this essentially is the condition that the probabilities must add up to one).
In general, the ensemble evolves over time according to the von Neumann equation.
Equilibrium ensembles (those that do not evolve over time, ) can be written solely as a function of conserved variables. For example, the microcanonical ensemble and canonical ensemble are strictly functions of the total energy, which is measured by the total energy operator (Hamiltonian). The grand canonical ensemble is additionally a function of the particle number, measured by the total particle number operator . Such equilibrium ensembles are a diagonal matrix in the orthogonal basis of states that simultaneously diagonalize each conserved variable. In bra–ket notation, the density matrix is
where the , indexed by , are the elements of a complete and orthogonal basis. (Note that in other bases, the density matrix is not necessarily diagonal.)
Classical mechanical
In classical mechanics, an ensemble is represented by a probability density function defined over the system's phase space. While an individual system evolves according to Hamilton's equations, the density function (the ensemble) evolves over time according to Liouville's equation.
In a mechanical system with a defined number of parts, the phase space has generalized coordinates called , and associated canonical momenta called . The ensemble is then represented by a joint probability density function .
If the number of parts in the system is allowed to vary among the systems in the ensemble (as in a grand ensemble where the number of particles is a random quantity), then it is a probability distribution over an extended phase space that includes further variables such as particle numbers (first kind of particle), (second kind of particle), and so on up to (the last kind of particle; is how many different kinds of particles there are). The ensemble is then represented by a joint probability density function . The number of coordinates varies with the numbers of particles.
Any mechanical quantity can be written as a function of the system's phase. The expectation value of any such quantity is given by an integral over the entire phase space of this quantity weighted by :
The condition of probability normalization applies, requiring
Phase space is a continuous space containing an infinite number of distinct physical states within any small region. In order to connect the probability density in phase space to a probability distribution over microstates, it is necessary to somehow partition the phase space into blocks that are distributed representing the different states of the system in a fair way. It turns out that the correct way to do this simply results in equal-sized blocks of canonical phase space, and so a microstate in classical mechanics is an extended region in the phase space of canonical coordinates that has a particular volume. In particular, the probability density function in phase space, , is related to the probability distribution over microstates, by a factor
where
is an arbitrary but predetermined constant with the units of , setting the extent of the microstate and providing correct dimensions to .
is an overcounting correction factor (see below), generally dependent on the number of particles and similar concerns.
Since can be chosen arbitrarily, the notional size of a microstate is also arbitrary. Still, the value of influences the offsets of quantities such as entropy and chemical potential, and so it is important to be consistent with the value of when comparing different systems.
Correcting overcounting in phase space
Typically, the phase space contains duplicates of the same physical state in multiple distinct locations. This is a consequence of the way that a physical state is encoded into mathematical coordinates; the simplest choice of coordinate system often allows a state to be encoded in multiple ways. An example of this is a gas of identical particles whose state is written in terms of the particles' individual positions and momenta: when two particles are exchanged, the resulting point in phase space is different, and yet it corresponds to an identical physical state of the system. It is important in statistical mechanics (a theory about physical states) to recognize that the phase space is just a mathematical construction, and to not naively overcount actual physical states when integrating over phase space. Overcounting can cause serious problems:
Dependence of derived quantities (such as entropy and chemical potential) on the choice of coordinate system, since one coordinate system might show more or less overcounting than another.
Erroneous conclusions that are inconsistent with physical experience, as in the mixing paradox.
Foundational issues in defining the chemical potential and the grand canonical ensemble.
It is in general difficult to find a coordinate system that uniquely encodes each physical state. As a result, it is usually necessary to use a coordinate system with multiple copies of each state, and then to recognize and remove the overcounting.
A crude way to remove the overcounting would be to manually define a subregion of phase space that includes each physical state only once and then exclude all other parts of phase space. In a gas, for example, one could include only those phases where the particles' coordinates are sorted in ascending order. While this would solve the problem, the resulting integral over phase space would be tedious to perform due to its unusual boundary shape. (In this case, the factor introduced above would be set to , and the integral would be restricted to the selected subregion of phase space.)
A simpler way to correct the overcounting is to integrate over all of phase space but to reduce the weight of each phase in order to exactly compensate the overcounting. This is accomplished by the factor introduced above, which is a whole number that represents how many ways a physical state can be represented in phase space. Its value does not vary with the continuous canonical coordinates, so overcounting can be corrected simply by integrating over the full range of canonical coordinates, then dividing the result by the overcounting factor. However, does vary strongly with discrete variables such as numbers of particles, and so it must be applied before summing over particle numbers.
As mentioned above, the classic example of this overcounting is for a fluid system containing various kinds of particles, where any two particles of the same kind are indistinguishable and exchangeable. When the state is written in terms of the particles' individual positions and momenta, then the overcounting related to the exchange of identical particles is corrected by using
This is known as "correct Boltzmann counting".
Ensembles in statistics
The formulation of statistical ensembles used in physics has now been widely adopted in other fields, in part because it has been recognized that the canonical ensemble or Gibbs measure serves to maximize the entropy of a system, subject to a set of constraints: this is the principle of maximum entropy. This principle has now been widely applied to problems in linguistics, robotics, and the like.
In addition, statistical ensembles in physics are often built on a principle of locality: that all interactions are only between neighboring atoms or nearby molecules. Thus, for example, lattice models, such as the Ising model, model ferromagnetic materials by means of nearest-neighbor interactions between spins. The statistical formulation of the principle of locality is now seen to be a form of the Markov property in the broad sense; nearest neighbors are now Markov blankets. Thus, the general notion of a statistical ensemble with nearest-neighbor interactions leads to Markov random fields, which again find broad applicability; for example in Hopfield networks.
Ensemble average
In statistical mechanics, the ensemble average is defined as the mean of a quantity that is a function of the microstate of a system, according to the distribution of the system on its micro-states in this ensemble.
Since the ensemble average is dependent on the ensemble chosen, its mathematical expression varies from ensemble to ensemble. However, the mean obtained for a given physical quantity does not depend on the ensemble chosen at the thermodynamic limit.
The grand canonical ensemble is an example of an open system.
Classical statistical mechanics
For a classical system in thermal equilibrium with its environment, the ensemble average takes the form of an integral over the phase space of the system:
where
is the ensemble average of the system property A,
is , known as thermodynamic beta,
H is the Hamiltonian of the classical system in terms of the set of coordinates and their conjugate generalized momenta ,
is the volume element of the classical phase space of interest.
The denominator in this expression is known as the partition function and is denoted by the letter Z.
Quantum statistical mechanics
In quantum statistical mechanics, for a quantum system in thermal equilibrium with its environment, the weighted average takes the form of a sum over quantum energy states, rather than a continuous integral:
Canonical ensemble average
The generalized version of the partition function provides the complete framework for working with ensemble averages in thermodynamics, information theory, statistical mechanics and quantum mechanics.
The microcanonical ensemble represents an isolated system in which energy (E), volume (V) and the number of particles (N) are all constant. The canonical ensemble represents a closed system which can exchange energy (E) with its surroundings (usually a heat bath), but the volume (V) and the number of particles (N) are all constant. The grand canonical ensemble represents an open system which can exchange energy (E) and particles (N) with its surroundings, but the volume (V) is kept constant.
Operational interpretation
In the discussion given so far, while rigorous, we have taken for granted that the notion of an ensemble is valid a priori, as is commonly done in physical context. What has not been shown is that the ensemble itself (not the consequent results) is a precisely defined object mathematically. For instance,
It is not clear where this very large set of systems exists (for example, is it a gas of particles inside a container?)
It is not clear how to physically generate an ensemble.
In this section, we attempt to partially answer this question.
Suppose we have a preparation procedure for a system in a physics lab: For example, the procedure might involve a physical apparatus and some protocols for manipulating the apparatus. As a result of this preparation procedure, some system is produced and maintained in isolation for some small period of time. By repeating this laboratory preparation procedure we obtain a sequence of systems X1, X2, ...,Xk, which in our mathematical idealization, we assume is an infinite sequence of systems. The systems are similar in that they were all produced in the same way. This infinite sequence is an ensemble.
In a laboratory setting, each one of these prepped systems might be used as input for one subsequent testing procedure. Again, the testing procedure involves a physical apparatus and some protocols; as a result of the testing procedure we obtain a yes or no answer. Given a testing procedure E applied to each prepared system, we obtain a sequence of values Meas (E, X1), Meas (E, X2), ..., Meas (E, Xk). Each one of these values is a 0 (or no) or a 1 (yes).
Assume the following time average exists:
For quantum mechanical systems, an important assumption made in the quantum logic approach to quantum mechanics is the identification of yes–no questions to the lattice of closed subspaces of a Hilbert space. With some additional technical assumptions one can then infer that states are given by density operators S so that:
We see this reflects the definition of quantum states in general: A quantum state is a mapping from the observables to their expectation values.
See also
Density matrix
Ensemble (fluid mechanics)
Phase space
Liouville's theorem (Hamiltonian)
Maxwell–Boltzmann statistics
Replication (statistics)
Notes
References
External links
Monte Carlo applet applied in statistical physics problems.
Equations of physics
Philosophy of thermal and statistical physics | Ensemble (mathematical physics) | [
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"Chemistry",
"Mathematics"
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59,211 | https://en.wikipedia.org/wiki/Altitude | Altitude is a distance measurement, usually in the vertical or "up" direction, between a reference datum and a point or object. The exact definition and reference datum varies according to the context (e.g., aviation, geometry, geographical survey, sport, or atmospheric pressure). Although the term altitude is commonly used to mean the height above sea level of a location, in geography the term elevation is often preferred for this usage.
In aviation, altitude is typically measured relative to mean sea level or above ground level to ensure safe navigation and flight operations. In geometry and geographical surveys, altitude helps create accurate topographic maps and understand the terrain's elevation. For high-altitude trekking and sports, knowing and adapting to altitude is vital for performance and safety. Higher altitudes mean reduced oxygen levels, which can lead to altitude sickness if proper acclimatization measures are not taken.
Vertical distance measurements in the "down" direction are commonly referred to as depth.
In aviation
The term altitude can have several meanings, and is always qualified by explicitly adding a modifier (e.g. "true altitude"), or implicitly through the context of the communication. Parties exchanging altitude information must be clear which definition is being used.
Aviation altitude is measured using either mean sea level (MSL) or local ground level (above ground level, or AGL) as the reference datum.
Pressure altitude divided by 100 feet (30 m) is the flight level, and is used above the transition altitude ( in the US, but may be as low as in other jurisdictions). So when the altimeter reads the country-specific flight level on the standard pressure setting the aircraft is said to be at "Flight level XXX/100" (where XXX is the transition altitude). When flying at a flight level, the altimeter is always set to standard pressure (29.92 inHg or 1013.25 hPa).
On the flight deck, the definitive instrument for measuring altitude is the pressure altimeter, which is an aneroid barometer with a front face indicating distance (feet or metres) instead of atmospheric pressure.
There are several types of altitude in aviation:
Indicated altitude is the reading on the altimeter when it is set to the local barometric pressure at mean sea level. In UK aviation radiotelephony usage, the vertical distance of a level, a point or an object considered as a point, measured from mean sea level; this is referred to over the radio as altitude.(see QNH)
Absolute altitude is the vertical distance of the aircraft above the terrain over which it is flying. It can be measured using a radar altimeter (or "absolute altimeter"). Also referred to as "radar height" or feet/metres above ground level (AGL).
True altitude is the actual elevation above mean sea level. It is indicated altitude corrected for non-standard temperature and pressure.
Height is the vertical distance above a reference point, commonly the terrain elevation. In UK aviation radiotelephony usage, the vertical distance of a level, a point or an object considered as a point, measured from a specified datum; this is referred to over the radio as height, where the specified datum is the airfield elevation (see QFE)
Pressure altitude is the elevation above a standard datum air-pressure plane (typically, 1013.25 millibars or 29.92" Hg). Pressure altitude is used to indicate "flight level" which is the standard for altitude reporting in the U.S. in Class A airspace (above roughly 18,000 feet). Pressure altitude and indicated altitude are the same when the altimeter setting is 29.92" Hg or 1013.25 millibars.
Density altitude is the altitude corrected for non-ISA International Standard Atmosphere atmospheric conditions. Aircraft performance depends on density altitude, which is affected by barometric pressure, humidity and temperature. On a very hot day, density altitude at an airport (especially one at a high elevation) may be so high as to preclude takeoff, particularly for helicopters or a heavily loaded aircraft.
These types of altitude can be explained more simply as various ways of measuring the altitude:
Indicated altitude – the altitude shown on the altimeter.
Absolute altitude – altitude in terms of the distance above the ground directly below
True altitude – altitude in terms of elevation above sea level
Height – vertical distance above a certain point
Pressure altitude – the air pressure in terms of altitude in the International Standard Atmosphere
Density altitude – the density of the air in terms of altitude in the International Standard Atmosphere in the air
In satellite orbits
In atmospheric studies
Atmospheric layers
The Earth's atmosphere is divided into several altitude regions. These regions start and finish at varying heights depending on season and distance from the poles. The altitudes stated below are averages:
Troposphere: surface to at the poles, at the Equator, ending at the Tropopause
Stratosphere: Troposphere to
Mesosphere: Stratosphere to
Thermosphere: Mesosphere to
Exosphere: Thermosphere to
The Kármán line, at an altitude of above sea level, by convention defines represents the demarcation between the atmosphere and space. The thermosphere and exosphere (along with the higher parts of the mesosphere) are regions of the atmosphere that are conventionally defined as space.
High altitude and low pressure
Regions on the Earth's surface (or in its atmosphere) that are high above mean sea level are referred to as high altitude. High altitude is sometimes defined to begin at above sea level.
At high altitude, atmospheric pressure is lower than that at sea level. This is due to two competing physical effects: gravity, which causes the air to be as close as possible to the ground; and the heat content of the air, which causes the molecules to bounce off each other and expand.
Temperature profile
The temperature profile of the atmosphere is a result of an interaction between radiation and convection. Sunlight in the visible spectrum hits the ground and heats it. The ground then heats the air at the surface. If radiation were the only way to transfer heat from the ground to space, the greenhouse effect of gases in the atmosphere would keep the ground at roughly , and the temperature would decay exponentially with height.
However, when air is hot, it tends to expand, which lowers its density. Thus, hot air tends to rise and transfer heat upward. This is the process of convection. Convection comes to equilibrium when a parcel of air at a given altitude has the same density as its surroundings. Air is a poor conductor of heat, so a parcel of air will rise and fall without exchanging heat. This is known as an adiabatic process, which has a characteristic pressure-temperature curve. As the pressure gets lower, the temperature decreases. The rate of decrease of temperature with elevation is known as the adiabatic lapse rate, which is approximately 9.8 °C per kilometer (or per 1000 feet) of altitude.
The presence of water in the atmosphere complicates the process of convection. Water vapor contains latent heat of vaporization. As air rises and cools, it eventually becomes saturated and cannot hold its quantity of water vapor. The water vapor condenses (forming clouds), and releases heat, which changes the lapse rate from the dry adiabatic lapse rate to the moist adiabatic lapse rate (5.5 °C per kilometer or per 1000 feet).
As an average, the International Civil Aviation Organization (ICAO) defines an international standard atmosphere (ISA) with a temperature lapse rate of 6.49 °C per kilometer (3.56 °F per 1,000 feet). The actual lapse rate can vary by altitude and by location.
Finally, only the troposphere (up to approximately of altitude) in the Earth's atmosphere undergoes notable convection; in the stratosphere, there is little vertical convection.
Effects on organisms
Humans
Medicine recognizes that altitudes above start to affect humans, and there is no record of humans living at extreme altitudes above for more than two years. As the altitude increases, atmospheric pressure decreases, which affects humans by reducing the partial pressure of oxygen. The lack of oxygen above can cause serious illnesses such as altitude sickness, high altitude pulmonary edema, and high altitude cerebral edema. The higher the altitude, the more likely are serious effects. The human body can adapt to high altitude by breathing faster, having a higher heart rate, and adjusting its blood chemistry. It can take days or weeks to adapt to high altitude. However, above , (in the "death zone"), altitude acclimatization becomes impossible.
There is a significantly lower overall mortality rate for permanent residents at higher altitudes. Additionally, there is a dose response relationship between increasing elevation and decreasing obesity prevalence in the United States. In addition, the recent hypothesis suggests that high altitude could be protective against Alzheimer's disease via action of erythropoietin, a hormone released by kidney in response to hypoxia.
However, people living at higher elevations have a statistically significant higher rate of suicide. The cause for the increased suicide risk is unknown so far.
Athletes
For athletes, high altitude produces two contradictory effects on performance. For explosive events (sprints up to 400 metres, long jump, triple jump) the reduction in atmospheric pressure signifies less atmospheric resistance, which generally results in improved athletic performance. For endurance events (races of 5,000 metres or more) the predominant effect is the reduction in oxygen which generally reduces the athlete's performance at high altitude. Sports organizations acknowledge the effects of altitude on performance: the International Association of Athletic Federations (IAAF), for example, marks record performances achieved at an altitude greater than with the letter "A".
Athletes also can take advantage of altitude acclimatization to increase their performance. The same changes that help the body cope with high altitude increase performance back at sea level. These changes are the basis of altitude training which forms an integral part of the training of athletes in a number of endurance sports including track and field, distance running, triathlon, cycling and swimming.
Other organisms
Decreased oxygen availability and decreased temperature make life at high altitude challenging. Despite these environmental conditions, many species have been successfully adapted at high altitudes. Animals have developed physiological adaptations to enhance oxygen uptake and delivery to tissues which can be used to sustain metabolism. The strategies used by animals to adapt to high altitude depend on their morphology and phylogeny. For example, small mammals face the challenge of maintaining body heat in cold temperatures, due to their small volume to surface area ratio. As oxygen is used as a source of metabolic heat production, the hypobaric hypoxia at high altitudes is problematic.
There is also a general trend of smaller body sizes and lower species richness at high altitudes, likely due to lower oxygen partial pressures. These factors may decrease productivity in high altitude habitats, meaning there will be less energy available for consumption, growth, and activity.
However, some species, such as birds, thrive at high altitude. Birds thrive because of physiological features that are advantageous for high-altitude flight.
See also
Atmosphere of Earth
Coffin corner (aerodynamics) At higher altitudes, the air density is lower than at sea level. At a certain altitude it is very difficult to keep an airplane in stable flight.
Geocentric altitude
Near space
References
External links
Downloadable ETOPO2 Raw Data Database (2 minute grid)
Downloadable ETOPO5 Raw Data Database (5 minute grid)
Aerospace
Physical geography
Topography
Vertical position | Altitude | [
"Physics"
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"Vertical position",
"Aerospace",
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59,220 | https://en.wikipedia.org/wiki/Base%20%28topology%29 | In mathematics, a base (or basis; : bases) for the topology of a topological space is a family of open subsets of such that every open set of the topology is equal to the union of some sub-family of . For example, the set of all open intervals in the real number line is a basis for the Euclidean topology on because every open interval is an open set, and also every open subset of can be written as a union of some family of open intervals.
Bases are ubiquitous throughout topology. The sets in a base for a topology, which are called , are often easier to describe and use than arbitrary open sets. Many important topological definitions such as continuity and convergence can be checked using only basic open sets instead of arbitrary open sets. Some topologies have a base of open sets with specific useful properties that may make checking such topological definitions easier.
Not all families of subsets of a set form a base for a topology on . Under some conditions detailed below, a family of subsets will form a base for a (unique) topology on , obtained by taking all possible unions of subfamilies. Such families of sets are very frequently used to define topologies. A weaker notion related to bases is that of a subbase for a topology. Bases for topologies are also closely related to neighborhood bases.
Definition and basic properties
Given a topological space , a base (or basis) for the topology (also called a base for if the topology is understood) is a family of open sets such that every open set of the topology can be represented as the union of some subfamily of . The elements of are called basic open sets.
Equivalently, a family of subsets of is a base for the topology if and only if and for every open set in and point there is some basic open set such that .
For example, the collection of all open intervals in the real line forms a base for the standard topology on the real numbers. More generally, in a metric space the collection of all open balls about points of forms a base for the topology.
In general, a topological space can have many bases. The whole topology is always a base for itself (that is, is a base for ). For the real line, the collection of all open intervals is a base for the topology. So is the collection of all open intervals with rational endpoints, or the collection of all open intervals with irrational endpoints, for example. Note that two different bases need not have any basic open set in common. One of the topological properties of a space is the minimum cardinality of a base for its topology, called the weight of and denoted . From the examples above, the real line has countable weight.
If is a base for the topology of a space , it satisfies the following properties:
(B1) The elements of cover , i.e., every point belongs to some element of .
(B2) For every and every point , there exists some such that .
Property (B1) corresponds to the fact that is an open set; property (B2) corresponds to the fact that is an open set.
Conversely, suppose is just a set without any topology and is a family of subsets of satisfying properties (B1) and (B2). Then is a base for the topology that it generates. More precisely, let be the family of all subsets of that are unions of subfamilies of Then is a topology on and is a base for .
(Sketch: defines a topology because it is stable under arbitrary unions by construction, it is stable under finite intersections by (B2), it contains by (B1), and it contains the empty set as the union of the empty subfamily of . The family is then a base for by construction.) Such families of sets are a very common way of defining a topology.
In general, if is a set and is an arbitrary collection of subsets of , there is a (unique) smallest topology on containing . (This topology is the intersection of all topologies on containing .) The topology is called the topology generated by , and is called a subbase for . The topology can also be characterized as the set of all arbitrary unions of finite intersections of elements of . (See the article about subbase.) Now, if also satisfies properties (B1) and (B2), the topology generated by can be described in a simpler way without having to take intersections: is the set of all unions of elements of (and is a base for in that case).
There is often an easy way to check condition (B2). If the intersection of any two elements of is itself an element of or is empty, then condition (B2) is automatically satisfied (by taking ). For example, the Euclidean topology on the plane admits as a base the set of all open rectangles with horizontal and vertical sides, and a nonempty intersection of two such basic open sets is also a basic open set. But another base for the same topology is the collection of all open disks; and here the full (B2) condition is necessary.
An example of a collection of open sets that is not a base is the set of all semi-infinite intervals of the forms and with . The topology generated by contains all open intervals , hence generates the standard topology on the real line. But is only a subbase for the topology, not a base: a finite open interval does not contain any element of (equivalently, property (B2) does not hold).
Examples
The set of all open intervals in forms a basis for the Euclidean topology on .
A non-empty family of subsets of a set that is closed under finite intersections of two or more sets, which is called a -system on , is necessarily a base for a topology on if and only if it covers . By definition, every σ-algebra, every filter (and so in particular, every neighborhood filter), and every topology is a covering -system and so also a base for a topology. In fact, if is a filter on then is a topology on and is a basis for it. A base for a topology does not have to be closed under finite intersections and many are not. But nevertheless, many topologies are defined by bases that are also closed under finite intersections. For example, each of the following families of subset of is closed under finite intersections and so each forms a basis for some topology on :
The set of all bounded open intervals in generates the usual Euclidean topology on .
The set of all bounded closed intervals in generates the discrete topology on and so the Euclidean topology is a subset of this topology. This is despite the fact that is not a subset of . Consequently, the topology generated by , which is the Euclidean topology on , is coarser than the topology generated by . In fact, it is strictly coarser because contains non-empty compact sets which are never open in the Euclidean topology.
The set of all intervals in such that both endpoints of the interval are rational numbers generates the same topology as . This remains true if each instance of the symbol is replaced by .
generates a topology that is strictly coarser than the topology generated by . No element of is open in the Euclidean topology on .
generates a topology that is strictly coarser than both the Euclidean topology and the topology generated by . The sets and are disjoint, but nevertheless is a subset of the topology generated by .
Objects defined in terms of bases
The order topology on a totally ordered set admits a collection of open-interval-like sets as a base.
In a metric space the collection of all open balls forms a base for the topology.
The discrete topology has the collection of all singletons as a base.
A second-countable space is one that has a countable base.
The Zariski topology on the spectrum of a ring has a base consisting of open sets that have specific useful properties. For the usual base for this topology, every finite intersection of basic open sets is a basic open set.
The Zariski topology of is the topology that has the algebraic sets as closed sets. It has a base formed by the set complements of algebraic hypersurfaces.
The Zariski topology of the spectrum of a ring (the set of the prime ideals) has a base such that each element consists of all prime ideals that do not contain a given element of the ring.
Theorems
A topology is finer than a topology if and only if for each and each basic open set of containing , there is a basic open set of containing and contained in .
If are bases for the topologies then the collection of all set products with each is a base for the product topology In the case of an infinite product, this still applies, except that all but finitely many of the base elements must be the entire space.
Let be a base for and let be a subspace of . Then if we intersect each element of with , the resulting collection of sets is a base for the subspace .
If a function maps every basic open set of into an open set of , it is an open map. Similarly, if every preimage of a basic open set of is open in , then is continuous.
is a base for a topological space if and only if the subcollection of elements of which contain form a local base at , for any point .
Base for the closed sets
Closed sets are equally adept at describing the topology of a space. There is, therefore, a dual notion of a base for the closed sets of a topological space. Given a topological space a family of closed sets forms a base for the closed sets if and only if for each closed set and each point not in there exists an element of containing but not containing
A family is a base for the closed sets of if and only if its in that is the family of complements of members of , is a base for the open sets of
Let be a base for the closed sets of Then
For each the union is the intersection of some subfamily of (that is, for any not in there is some containing and not containing ).
Any collection of subsets of a set satisfying these properties forms a base for the closed sets of a topology on The closed sets of this topology are precisely the intersections of members of
In some cases it is more convenient to use a base for the closed sets rather than the open ones. For example, a space is completely regular if and only if the zero sets form a base for the closed sets. Given any topological space the zero sets form the base for the closed sets of some topology on This topology will be the finest completely regular topology on coarser than the original one. In a similar vein, the Zariski topology on An is defined by taking the zero sets of polynomial functions as a base for the closed sets.
Weight and character
We shall work with notions established in .
Fix X a topological space. Here, a network is a family of sets, for which, for all points x and open neighbourhoods U containing x, there exists B in for which Note that, unlike a basis, the sets in a network need not be open.
We define the weight, w(X), as the minimum cardinality of a basis; we define the network weight, nw(X), as the minimum cardinality of a network; the character of a point, as the minimum cardinality of a neighbourhood basis for x in X; and the character of X to be
The point of computing the character and weight is to be able to tell what sort of bases and local bases can exist. We have the following facts:
nw(X) ≤ w(X).
if X is discrete, then w(X) = nw(X) = |X|.
if X is Hausdorff, then nw(X) is finite if and only if X is finite discrete.
if B is a basis of X then there is a basis of size
if N a neighbourhood basis for x in X then there is a neighbourhood basis of size
if is a continuous surjection, then nw(Y) ≤ w(X). (Simply consider the Y-network for each basis B of X.)
if is Hausdorff, then there exists a weaker Hausdorff topology so that So a fortiori, if X is also compact, then such topologies coincide and hence we have, combined with the first fact, nw(X) = w(X).
if a continuous surjective map from a compact metrizable space to an Hausdorff space, then Y is compact metrizable.
The last fact follows from f(X) being compact Hausdorff, and hence (since compact metrizable spaces are necessarily second countable); as well as the fact that compact Hausdorff spaces are metrizable exactly in case they are second countable. (An application of this, for instance, is that every path in a Hausdorff space is compact metrizable.)
Increasing chains of open sets
Using the above notation, suppose that w(X) ≤ κ some infinite cardinal. Then there does not exist a strictly increasing sequence of open sets (equivalently strictly decreasing sequence of closed sets) of length ≥ κ+.
To see this (without the axiom of choice), fix
as a basis of open sets. And suppose per contra, that
were a strictly increasing sequence of open sets. This means
For
we may use the basis to find some Uγ with x in Uγ ⊆ Vα. In this way we may well-define a map, f : κ+ → κ mapping each α to the least γ for which Uγ ⊆ Vα and meets
This map is injective, otherwise there would be α < β with f(α) = f(β) = γ, which would further imply Uγ ⊆ Vα but also meets
which is a contradiction. But this would go to show that κ+ ≤ κ, a contradiction.
See also
Esenin-Volpin's theorem
Gluing axiom
Neighbourhood system
Notes
References
Bibliography
General topology | Base (topology) | [
"Mathematics"
] | 2,856 | [
"General topology",
"Topology"
] |
59,358 | https://en.wikipedia.org/wiki/Mycorrhiza | A mycorrhiza (; , mycorrhiza, or mycorrhizas) is a symbiotic association between a fungus and a plant. The term mycorrhiza refers to the role of the fungus in the plant's rhizosphere, the plant root system and its surroundings. Mycorrhizae play important roles in plant nutrition, soil biology, and soil chemistry.
In a mycorrhizal association, the fungus colonizes the host plant's root tissues, either intracellularly as in arbuscular mycorrhizal fungi, or extracellularly as in ectomycorrhizal fungi. The association is normally mutualistic. In particular species, or in particular circumstances, mycorrhizae may have a parasitic association with host plants.
Definition
A mycorrhiza is a symbiotic association between a green plant and a fungus. The plant makes organic molecules by photosynthesis and supplies them to the fungus in the form of sugars or lipids, while the fungus supplies the plant with water and mineral nutrients, such as phosphorus, taken from the soil. Mycorrhizas are located in the roots of vascular plants, but mycorrhiza-like associations also occur in bryophytes and there is fossil evidence that early land plants that lacked roots formed arbuscular mycorrhizal associations. Most plant species form mycorrhizal associations, though some families like Brassicaceae and Chenopodiaceae cannot. Different forms for the association are detailed in the next section. The most common is the arbuscular type that is present in 70% of plant species, including many crop plants such as cereals and legumes.
Evolution
Fossil and genetic evidence indicate that mycorrhizae are ancient, potentially as old as the terrestrialization of plants. Genetic evidence indicates that all land plants share a single common ancestor, which appears to have quickly adopted mycorrhizal symbiosis, and research suggests that proto-mycorrhizal fungi were a key factor enabling plant terrestrialization. The 400 million year old Rhynie chert contains an assemblage of fossil plants preserved in sufficient detail that arbuscular mycorrhizae have been observed in the stems of Aglaophyton major, giving a lower bound for how late mycorrhizal symbiosis may have developed. Ectomycorrhizae developed substantially later, during the Jurassic period, while most other modern mycorrhizal families, including orchid and ericoid mycorrhizae, date to the period of angiosperm radiation in the Cretaceous period. There is genetic evidence that the symbiosis between legumes and nitrogen-fixing bacteria is an extension of mycorrhizal symbiosis. The modern distribution of mycorrhizal fungi appears to reflect an increasing complexity and competition in root morphology associated with the dominance of angiosperms in the Cenozoic Era, characterized by complex ecological dynamics between species.
Types
The mycorrhizal lifestyle has independently convergently evolved multiple times in the history of Earth. There are multiple ways to categorize mycorrhizal symbiosis. One major categorization is the division between ectomycorrhizas and endomycorrhizas. The two types are differentiated by the fact that the hyphae of ectomycorrhizal fungi do not penetrate individual cells within the root, while the hyphae of endomycorrhizal fungi penetrate the cell wall and invaginate the cell membrane.
Similar symbiotic relationships
Some forms of plant-fungal symbiosis are similar to mycorrhizae, but considered distinct. One example is fungal endophytes. Endophytes are defined as organisms that can live within plant cells without causing harm to the plant. They are distinguishable from mycorrhizal fungi by the absence of nutrient-transferring structures for bringing in nutrients from outside the plant. Some lineages of mycorrhizal fungi may have evolved from endophytes into mycorrhizal fungi, and some fungi can live as mycorrhizae or as endophytes.
Ectomycorrhiza
Ectomycorrhizae are distinct in that they do not penetrate into plant cells, but instead form a structure called a Hartig net that penetrates between cells. Ectomycorrhizas consist of a hyphal sheath, or mantle, covering the root tip and the Hartig net of hyphae surrounding the plant cells within the root cortex. In some cases the hyphae may also penetrate the plant cells, in which case the mycorrhiza is called an endomycorrhiza. Outside the root, ectomycorrhizal extramatrical mycelium forms an extensive network within the soil and leaf litter. Other forms of mycorrhizae, including arbuscular, ericoid, arbutoid, monotropoid, and orchid mycorrhizas, are considered endomycorrhizae.
Ectomycorrhizas, or EcM, are symbiotic associations between the roots of around 10% of plant families, mostly woody plants including the birch, dipterocarp, eucalyptus, oak, pine, and rose families, orchids, and fungi belonging to the Basidiomycota, Ascomycota, and Zygomycota. Ectomycorrhizae associate with relatively few plant species, only about 2% of plant species on Earth, but the species they associate with are mostly trees and woody plants that are highly dominant in their ecosystems, meaning plants in ectomycorrhizal relationships make up a large proportion of plant biomass. Some EcM fungi, such as many Leccinum and Suillus, are symbiotic with only one particular genus of plant, while other fungi, such as the Amanita, are generalists that form mycorrhizas with many different plants. An individual tree may have 15 or more different fungal EcM partners at one time. While the diversity of plants involved in EcM is low, the diversity of fungi involved in EcM is high. Thousands of ectomycorrhizal fungal species exist, hosted in over 200 genera. A recent study has conservatively estimated global ectomycorrhizal fungal species richness at approximately 7750 species, although, on the basis of estimates of knowns and unknowns in macromycete diversity, a final estimate of ECM species richness would probably be between 20,000 and 25,000. Ectomycorrhizal fungi evolved independently from saprotrophic ancestors many times in the group's history.
Nutrients can be shown to move between different plants through the fungal network. Carbon has been shown to move from paper birch seedlings into adjacent Douglas-fir seedlings, although not conclusively through a common mycorrhizal network, thereby promoting succession in ecosystems. The ectomycorrhizal fungus Laccaria bicolor has been found to lure and kill springtails to obtain nitrogen, some of which may then be transferred to the mycorrhizal host plant. In a study by Klironomos and Hart, Eastern White Pine inoculated with L. bicolor was able to derive up to 25% of its nitrogen from springtails. When compared with non-mycorrhizal fine roots, ectomycorrhizae may contain very high concentrations of trace elements, including toxic metals (cadmium, silver) or chlorine.
The first genomic sequence for a representative of symbiotic fungi, the ectomycorrhizal basidiomycete L. bicolor, was published in 2008. An expansion of several multigene families occurred in this fungus, suggesting that adaptation to symbiosis proceeded by gene duplication. Within lineage-specific genes those coding for symbiosis-regulated secreted proteins showed an up-regulated expression in ectomycorrhizal root tips suggesting a role in the partner communication. L. bicolor is lacking enzymes involved in the degradation of plant cell wall components (cellulose, hemicellulose, pectins and pectates), preventing the symbiont from degrading host cells during the root colonisation. By contrast, L. bicolor possesses expanded multigene families associated with hydrolysis of bacterial and microfauna polysaccharides and proteins. This genome analysis revealed the dual saprotrophic and biotrophic lifestyle of the mycorrhizal fungus that enables it to grow within both soil and living plant roots. Since then, the genomes of many other ectomycorrhizal fungal species have been sequenced further expanding the study of gene families and evolution in these organisms.
Arbutoid mycorrhiza
This type of mycorrhiza involves plants of the Ericaceae subfamily Arbutoideae. It is however different from ericoid mycorrhiza and resembles ectomycorrhiza, both functionally and in terms of the fungi involved. It differs from ectomycorrhiza in that some hyphae actually penetrate into the root cells, making this type of mycorrhiza an ectendomycorrhiza.
Arbuscular mycorrhiza
Arbuscular mycorrhizas, (formerly known as vesicular-arbuscular mycorrhizas), have hyphae that penetrate plant cells, producing branching, tree-like structures called arbuscules within the plant cells for nutrient exchange. Often, balloon-like storage structures, termed vesicles, are also produced. In this interaction, fungal hyphae do not in fact penetrate the protoplast (i.e. the interior of the cell), but invaginate the cell membrane, creating a so-called peri-arbuscular membrane. The structure of the arbuscules greatly increases the contact surface area between the hypha and the host cell cytoplasm to facilitate the transfer of nutrients between them. Arbuscular mycorrhizas are obligate biotrophs, meaning that they depend upon the plant host for both growth and reproduction; they have lost the ability to sustain themselves by decomposing dead plant material. Twenty percent of the photosynthetic products made by the plant host are consumed by the fungi, the transfer of carbon from the terrestrial host plant is then exchanged by equal amounts of phosphate from the fungi to the plant host.
Contrasting with the pattern seen in ectomycorrhizae, the species diversity of AMFs is very low, but the diversity of plant hosts is very high; an estimated 78% of all plant species associate with AMFs. Arbuscular mycorrhizas are formed only by fungi in the division Glomeromycota. Fossil evidence and DNA sequence analysis suggest that this mutualism appeared 400-460 million years ago, when the first plants were colonizing land. Arbuscular mycorrhizas are found in 85% of all plant families, and occur in many crop species. The hyphae of arbuscular mycorrhizal fungi produce the glycoprotein glomalin, which may be one of the major stores of carbon in the soil. Arbuscular mycorrhizal fungi have (possibly) been asexual for many millions of years and, unusually, individuals can contain many genetically different nuclei (a phenomenon called heterokaryosis).
Mucoromycotina fine root endophytes
Mycorrhizal fungi belonging to Mucoromycotina, known as “fine root endophytes" (MFREs), were mistakenly identified as arbuscular mycorrhizal fungi until recently. While similar to AMF, MFREs are from subphylum Mucoromycotina instead of Glomeromycotina. Their morphology when colonizing a plant root is very similar to AMF, but they form fine textured hyphae. Effects of MFREs may have been mistakenly attributed to AMFs due to confusion between the two, complicated by the fact that AMFs and MFREs often colonize the same hosts simultaneously. Unlike AMFs, they appear capable of surviving without a host. This group of mycorrhizal fungi is little understood, but appears to prefer wet, acidic soils and forms symbiotic relationships with liverworts, hornworts, lycophytes, and angiosperms.
Ericoid mycorrhiza
Ericoid mycorrhizae, or ErMs, involve only plants in Ericales and are the most recently evolved of the major mycorrhizal relationships. Plants that form ericoid mycorrhizae are mostly woody understory shrubs; hosts include blueberries, bilberries, cranberries, mountain laurels, rhododendrons, heather, neinei, and giant grass tree. ErMs are most common in boreal forests, but are found in two-thirds of all forests on Earth. Ericoid mycorrhizal fungi belong to several different lineages of fungi. Some species can live as endophytes entirely within plant cells even within plants outside the Ericales, or live independently as saprotrophs that decompose dead organic matter. This ability to switch between multiple lifestyle types makes ericoid mycorrhizal fungi very adaptable.
Plants that participate in these symbioses have specialized roots with no root hairs, which are covered with a layer of epidermal cells that the fungus penetrates into and completely occupies. The fungi have a simple intraradical (growth in cells) phase, consisting of dense coils of hyphae in the outermost layer of root cells. There is no periradical phase and the extraradical phase consists of sparse hyphae that don't extend very far into the surrounding soil. They might form sporocarps (probably in the form of small cups), but their reproductive biology is poorly understood.
Plants participating in ericoid mycorrhizal symbioses are found in acidic, nutrient-poor conditions. Whereas AMFs have lost their saprotrophic capabilities, and EcM fungi have significant variation in their ability to produce enzymes needed for a saprotrophic lifestyle, fungi involved in ErMs have fully retained the ability to decompose plant material for sustenance. Some ericoid mycorrhizal fungi have actually expanded their repertoire of enzymes for breaking down organic matter. They can extract nitrogen from cellulose, hemicellulose, lignin, pectin, and chitin. This would increase the benefit they can provide to their plant symbiotic partners.
Orchid mycorrhiza
All orchids are myco-heterotrophic at some stage during their lifecycle, meaning that they can survive only if they form orchid mycorrhizae. Orchid seeds are so small that they contain no nutrition to sustain the germinating seedling, and instead must gain the energy to grow from their fungal symbiont. The OM relationship is asymmetric; the plant seems to benefit more than the fungus, and some orchids are entirely mycoheterotrophic, lacking chlorophyll for photosynthesis. It is actually unknown whether fully autotrophic orchids that do not receive some of their carbon from fungi exist or not. Like fungi that form ErMs, OM fungi can sometimes live as endophytes or as independent saprotrophs. In the OM symbiosis, hyphae penetrate into the root cells and form pelotons (coils) for nutrient exchange.
Monotropoid mycorrhiza
This type of mycorrhiza occurs in the subfamily Monotropoideae of the Ericaceae, as well as several genera in the Orchidaceae. These plants are heterotrophic or mixotrophic and derive their carbon from the fungus partner. This is thus a non-mutualistic, parasitic type of mycorrhizal symbiosis.
Function
Mycorrhizal fungi form a mutualistic relationship with the roots of most plant species. In such a relationship, both the plants themselves and those parts of the roots that host the fungi, are said to be mycorrhizal. Relatively few of the mycorrhizal relationships between plant species and fungi have been examined to date, but 95% of the plant families investigated are predominantly mycorrhizal either in the sense that most of their species associate beneficially with mycorrhizae, or are absolutely dependent on mycorrhizae. The Orchidaceae are notorious as a family in which the absence of the correct mycorrhizae is fatal even to germinating seeds.
Recent research into ectomycorrhizal plants in boreal forests has indicated that mycorrhizal fungi and plants have a relationship that may be more complex than simply mutualistic. This relationship was noted when mycorrhizal fungi were unexpectedly found to be hoarding nitrogen from plant roots in times of nitrogen scarcity. Researchers argue that some mycorrhizae distribute nutrients based upon the environment with surrounding plants and other mycorrhizae. They go on to explain how this updated model could explain why mycorrhizae do not alleviate plant nitrogen limitation, and why plants can switch abruptly from a mixed strategy with both mycorrhizal and nonmycorrhizal roots to a purely mycorrhizal strategy as soil nitrogen availability declines. It has also been suggested that evolutionary and phylogenetic relationships can explain much more variation in the strength of mycorrhizal mutualisms than ecological factors.
Formation
To successfully engage in mutualistic symbiotic relationships with other organisms, such as mycorrhizal fungi and any of the thousands of microbes that colonize plants, plants must discriminate between mutualists and pathogens, allowing the mutualists to colonize while activating an immune response towards the pathogens. Plant genomes code for potentially hundreds of receptors for detecting chemical signals from other organisms. Plants dynamically adjust their symbiotic and immune responses, changing their interactions with their symbionts in response to feedbacks detected by the plant. In plants, the mycorrhizal symbiosis is regulated by the common symbiosis signaling pathway (CSSP), a set of genes involved in initiating and maintaining colonization by endosymbiotic fungi and other endosymbionts such as Rhizobia in legumes. The CSSP has origins predating the colonization of land by plants, demonstrating that the co-evolution of plants and arbuscular mycorrhizal fungi is over 500 million years old. In arbuscular mycorrhizal fungi, the presence of strigolactones, a plant hormone, secreted from roots induces fungal spores in the soil to germinate, stimulates their metabolism, growth and branching, and prompts the fungi to release chemical signals the plant can detect. Once the plant and fungus recognize one another as suitable symbionts, the plant activates the common symbiotic signaling pathway, which causes changes in the root tissues that enable the fungus to colonize.
Experiments with arbuscular mycorrhizal fungi have identified numerous chemical compounds to be involved in the "chemical dialog" that occurs between the prospective symbionts before symbiosis is begun. In plants, almost all plant hormones play a role in initiating or regulating AMF symbiosis, and other chemical compounds are also suspected to have a signaling function. While the signals emitted by the fungi are less understood, it has been shown that chitinaceous molecules known as Myc factors are essential for the formation of arbuscular mycorrhizae. Signals from plants are detected by LysM-containing receptor-like kinases, or LysM-RLKs. AMF genomes also code for potentially hundreds of effector proteins, of which only a few have a proven effect on mycorrhizal symbiosis, but many others likely have a function in communication with plant hosts as well.
Many factors are involved in the initiation of mycorrhizal symbiosis, but particularly influential is the plant's need for phosphorus. Experiments involving rice plants with a mutation disabling their ability to detect P starvation show that arbuscular mycorrhizal fungi detection, recruitment and colonization is prompted when the plant detects that it is starved of phosphorus. Nitrogen starvation also plays a role in initiating AMF symbiosis.
Mechanisms
The mechanisms by which mycorrhizae increase absorption include some that are physical and some that are chemical. Physically, most mycorrhizal mycelia are much smaller in diameter than the smallest root or root hair, and thus can explore soil material that roots and root hairs cannot reach, and provide a larger surface area for absorption. Chemically, the cell membrane chemistry of fungi differs from that of plants. For example, they may secrete organic acids that dissolve or chelate many ions, or release them from minerals by ion exchange. Mycorrhizae are especially beneficial for the plant partner in nutrient-poor soils.
Sugar-water/mineral exchange
The mycorrhizal mutualistic association provides the fungus with relatively constant and direct access to carbohydrates, such as glucose and sucrose. The carbohydrates are translocated from their source (usually leaves) to root tissue and on to the plant's fungal partners. In return, the plant gains the benefits of the mycelium's higher absorptive capacity for water and mineral nutrients, partly because of the large surface area of fungal hyphae, which are much longer and finer than plant root hairs, and partly because some such fungi can mobilize soil minerals unavailable to the plants' roots. The effect is thus to improve the plant's mineral absorption capabilities.
Unaided plant roots may be unable to take up nutrients that are chemically or physically immobilised; examples include phosphate ions and micronutrients such as iron. One form of such immobilization occurs in soil with high clay content, or soils with a strongly basic pH. The mycelium of the mycorrhizal fungus can, however, access many such nutrient sources, and make them available to the plants they colonize. Thus, many plants are able to obtain phosphate without using soil as a source. Another form of immobilisation is when nutrients are locked up in organic matter that is slow to decay, such as wood, and some mycorrhizal fungi act directly as decay organisms, mobilising the nutrients and passing some onto the host plants; for example, in some dystrophic forests, large amounts of phosphate and other nutrients are taken up by mycorrhizal hyphae acting directly on leaf litter, bypassing the need for soil uptake. Inga alley cropping, an agroforestry technique proposed as an alternative to slash and burn rainforest destruction, relies upon mycorrhiza within the root system of species of Inga to prevent the rain from washing phosphorus out of the soil.
In some more complex relationships, mycorrhizal fungi do not just collect immobilised soil nutrients, but connect individual plants together by mycorrhizal networks that transport water, carbon, and other nutrients directly from plant to plant through underground hyphal networks.
Suillus tomentosus, a basidiomycete fungus, produces specialized structures known as tuberculate ectomycorrhizae with its plant host lodgepole pine (Pinus contorta var. latifolia). These structures have been shown to host nitrogen fixing bacteria which contribute a significant amount of nitrogen and allow the pines to colonize nutrient-poor sites.
Disease, drought and salinity resistance and its correlation to mycorrhizae
Mycorrhizal plants are often more resistant to diseases, such as those caused by microbial soil-borne pathogens. These associations have been found to assist in plant defense both above and belowground. Mycorrhizas have been found to excrete enzymes that are toxic to soil borne organisms such as nematodes. More recent studies have shown that mycorrhizal associations result in a priming effect of plants that essentially acts as a primary immune response. When this association is formed a defense response is activated similarly to the response that occurs when the plant is under attack. As a result of this inoculation, defense responses are stronger in plants with mycorrhizal associations.
Ecosystem services provided by mycorrhizal fungi may depend on the soil microbiome. Furthermore, mycorrhizal fungi was significantly correlated with soil physical variable, but only with water level and not with aggregate stability and can lead also to more resistant to the effects of drought. Moreover, the significance of mycorrhizal fungi also includes alleviation of salt stress and its beneficial effects on plant growth and productivity. Although salinity can negatively affect mycorrhizal fungi, many reports show improved growth and performance of mycorrhizal plants under salt stress conditions.
Resistance to insects
Plants connected by mycorrhizal fungi in mycorrhizal networks can use these underground connections to communicate warning signals. For example, when a host plant is attacked by an aphid, the plant signals surrounding connected plants of its condition. Both the host plant and those connected to it release volatile organic compounds that repel aphids and attract parasitoid wasps, predators of aphids. This assists the mycorrhizal fungi by conserving its food supply.
Colonization of barren soil
Plants grown in sterile soils and growth media often perform poorly without the addition of spores or hyphae of mycorrhizal fungi to colonise the plant roots and aid in the uptake of soil mineral nutrients. The absence of mycorrhizal fungi can also slow plant growth in early succession or on degraded landscapes. The introduction of alien mycorrhizal plants to nutrient-deficient ecosystems puts indigenous non-mycorrhizal plants at a competitive disadvantage. This aptitude to colonize barren soil is defined by the category Oligotroph.
Resistance to toxicity
Fungi have a protective role for plants rooted in soils with high metal concentrations, such as acidic and contaminated soils. Pine trees inoculated with Pisolithus tinctorius planted in several contaminated sites displayed high tolerance to the prevailing contaminant, survivorship and growth. One study discovered the existence of Suillus luteus strains with varying tolerance of zinc. Another study discovered that zinc-tolerant strains of Suillus bovinus conferred resistance to plants of Pinus sylvestris. This was probably due to binding of the metal to the extramatricial mycelium of the fungus, without affecting the exchange of beneficial substances.
Occurrence of mycorrhizal associations
Mycorrhizas are present in 92% of plant families studied (80% of species), with arbuscular mycorrhizas being the ancestral and predominant form, and the most prevalent symbiotic association found in the plant kingdom. The structure of arbuscular mycorrhizas has been highly conserved since their first appearance in the fossil record, with both the development of ectomycorrhizas and the loss of mycorrhizas, evolving convergently on multiple occasions.
Associations of fungi with the roots of plants have been known since at least the mid-19th century. However, early observers simply recorded the fact without investigating the relationships between the two organisms. This symbiosis was studied and described by Franciszek Kamieński in 1879–1882.
Climate change
CO2 released by human activities is causing climate change and possible damage to mycorrhizae, but the direct effect of an increase in the gas should be to benefit plants and mycorrhizae. In Arctic regions, nitrogen and water are harder for plants to obtain, making mycorrhizae crucial to plant growth. Since mycorrhizae tend to do better in cooler temperatures, warming could be detrimental to them. Gases such as SO2, NO-x, and O3 produced by human activity may harm mycorrhizae, causing reduction in "propagules, the colonization of roots, degradation in connections between trees, reduction in the mycorrhizal incidence in trees, and reduction in the enzyme activity of ectomycorrhizal roots."
A company in Israel, Groundwork BioAg, has discovered a method of using mycorrhizal fungi to increase agricultural crops while sequestering greenhouse gases and eliminating CO2 from the atmosphere.
Conservation and mapping
In 2021, the Society for the Protection of Underground Networks was launched. SPUN is a science-based initiative to map and protect the mycorrhizal networks regulating Earth’s climate and ecosystems. Its stated goals are mapping, protecting, and harnessing mycorrhizal fungi.
See also
Effect of climate change on plant biodiversity
Endosymbiont
Epibiont, an organism that grows on another life form
Endophyte
Epiphyte
Epiphytic fungus
Mucigel
Mycorrhizal fungi and soil carbon storage
Mycorrhizal network
Rhizobia
Suzanne Simard
References
External links
International Mycorrhiza Society International Mycorrhiza Society
Mohamed Hijri: A simple solution to the coming phosphorus crisis video recommending agricultural mycorrhiza use to conserve phosphorus reserves & 85% waste problem @Ted.com
Mycorrhizal Associations: The Web Resource Comprehensive illustrations and lists of mycorrhizal and nonmycorrhizal plants and fungi
Mycorrhizas – a successful symbiosis Biosafety research into genetically modified barley
MycorWiki a portal concerned with the biology and ecology of ectomycorrhizal fungi and other forest fungi.
Plant roots
Soil biology
Symbiosis
Oligotrophs
Fungus ecology | Mycorrhiza | [
"Biology"
] | 6,259 | [
"Fungi",
"Symbiosis",
"Behavior",
"Biological interactions",
"Fungus ecology",
"Soil biology"
] |
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